Solutions to Axler’s Measure, Integration & Real Analysis

Solutions to every exercise in Sheldon Axler’s Measure, Integration & Real Analysis (Springer GTM 282, 2020) — 587 exercises across sections 1A–12. The book is open access at measure.axler.net and is filed at Measure, Integration & Real Analysis (Axler).

    /

Riemann Integration

Exercises 1A

Problem (1A.1)

Suppose \(f : [a, b] \to \mathbb{R}\) is a bounded function such that

\begin{equation*} L(f, P, [a, b]) = U(f, P, [a, b]) \end{equation*}

for some partition \(P\) of \([a, b]\). Prove that \(f\) is a constant function on \([a, b]\).

Solution

Let \(P\) be the partition \(a = x_0 < x_1 < \cdots < x_n = b\). By the definition of the lower and upper Riemann sums (1.3),

\begin{equation*} U(f, P, [a,b]) - L(f, P, [a,b]) = \sum_{j=1}^{n} (x_j - x_{j-1}) \Big( \sup_{[x_{j-1}, x_j]} f - \inf_{[x_{j-1}, x_j]} f \Big). \end{equation*}

For each \(j\) the number \(x_j - x_{j-1}\) is strictly positive, and the infimum of a nonempty set of real numbers is less than or equal to its supremum, so every summand above is nonnegative. The hypothesis says the whole sum equals \(0\). A finite sum of nonnegative numbers equals \(0\) only if each summand equals \(0\), and since \(x_j - x_{j-1} > 0\) we conclude

\begin{equation*} \sup_{[x_{j-1}, x_j]} f = \inf_{[x_{j-1}, x_j]} f \quad \text{for each } j = 1, \ldots, n. \end{equation*}

Fix \(j\). For every \(x \in [x_{j-1}, x_j]\) we have \(\inf_{[x_{j-1},x_j]} f \le f(x) \le \sup_{[x_{j-1},x_j]} f\), and the two outer quantities are equal; hence \(f\) is constant on \([x_{j-1}, x_j]\). Write \(c_j\) for that constant value.

It remains to see that \(c_1 = c_2 = \cdots = c_n\). For \(1 \le j \le n-1\) the point \(x_j\) belongs to both \([x_{j-1}, x_j]\) and \([x_j, x_{j+1}]\), so

\begin{equation*} c_j = f(x_j) = c_{j+1}. \end{equation*}

Thus all the \(c_j\) are equal to a common value \(c\). Since \([a,b] = [x_0, x_1] \cup [x_1, x_2] \cup \cdots \cup [x_{n-1}, x_n]\), every \(x \in [a,b]\) lies in some \([x_{j-1}, x_j]\) and therefore satisfies \(f(x) = c_j = c\). Hence \(f\) is constant on \([a,b]\).

Problem (1A.2)

Suppose \(a \le s < t \le b\). Define \(f : [a, b] \to \mathbb{R}\) by

\begin{equation*} f(x) = \begin{cases} 1 & \text{if } s < x < t, \\ 0 & \text{otherwise.} \end{cases} \end{equation*}

Prove that \(f\) is Riemann integrable on \([a, b]\) and that \(\int_a^b f = t - s\).

Solution

The function \(f\) takes only the values \(0\) and \(1\), so it is bounded and the lower and upper Riemann sums are defined.

Fix \(\varepsilon\) with \(0 < \varepsilon < \frac{t-s}{2}\), so that

\begin{equation*} s < s + \varepsilon < t - \varepsilon < t . \end{equation*}

Let \(P_\varepsilon\) be the partition of \([a,b]\) whose points are the distinct elements of

\begin{equation*} \{a,\; s,\; s+\varepsilon,\; t-\varepsilon,\; t,\; b\}, \end{equation*}

listed in increasing order. (All six numbers lie in \([a,b]\) and are listed in nondecreasing order; the only possible coincidences are \(s = a\) and \(t = b\), and deleting repetitions leaves a genuine partition of \([a,b]\).) Its subintervals are

\begin{equation*} [a,s] \ (\text{present only if } a < s), \quad [s, s+\varepsilon], \quad [s+\varepsilon, t-\varepsilon], \quad [t-\varepsilon, t], \quad [t,b] \ (\text{present only if } t < b). \end{equation*}

Lower sum. Every \(x \in [s+\varepsilon, t-\varepsilon]\) satisfies \(s < x < t\), so \(f \equiv 1\) there and \(\inf_{[s+\varepsilon,\,t-\varepsilon]} f = 1\). On every other subinterval the infimum of \(f\) is at least \(0\), since \(f \ge 0\). Hence

\begin{equation*} L(f, P_\varepsilon, [a,b]) \ge (t - \varepsilon) - (s + \varepsilon) = (t-s) - 2\varepsilon . \end{equation*}

Therefore \(L(f, [a,b]) \ge (t-s) - 2\varepsilon\) for every \(\varepsilon \in \big(0, \tfrac{t-s}{2}\big)\), and letting \(\varepsilon \to 0\) gives

\begin{equation*} L(f, [a,b]) \ge t - s . \end{equation*}

Upper sum. If \(x \in [a,s]\) then \(x \le s\), so \(f(x) = 0\); thus \(\sup_{[a,s]} f = 0\). If \(x \in [t,b]\) then \(x \ge t\), so \(f(x) = 0\); thus \(\sup_{[t,b]} f = 0\). On the remaining three subintervals we use only \(\sup f \le 1\). Hence

\begin{equation*} U(f, P_\varepsilon, [a,b]) \le 0 \cdot (s-a) + 1 \cdot \varepsilon + 1 \cdot \big((t-\varepsilon)-(s+\varepsilon)\big) + 1 \cdot \varepsilon + 0 \cdot (b-t) = t - s . \end{equation*}

Therefore \(U(f, [a,b]) \le t-s\).

Combining the two displays with 1.8 gives

\begin{equation*} t - s \le L(f, [a,b]) \le U(f, [a,b]) \le t - s , \end{equation*}

so all these quantities are equal. Thus \(f\) is Riemann integrable on \([a,b]\) and \(\int_a^b f = t - s\).

Problem (1A.3)

Suppose \(f : [a, b] \to \mathbb{R}\) is a bounded function. Prove that \(f\) is Riemann integrable if and only if for each \(\varepsilon > 0\), there exists a partition \(P\) of \([a, b]\) such that

\begin{equation*} U(f, P, [a, b]) - L(f, P, [a, b]) < \varepsilon . \end{equation*}

Solution

Throughout, recall from 1.7 that \(L(f,[a,b]) = \sup_P L(f,P,[a,b])\) and \(U(f,[a,b]) = \inf_P U(f,P,[a,b])\), so that

\begin{equation*} L(f, P, [a,b]) \le L(f, [a,b]) \quad\text{and}\quad U(f,[a,b]) \le U(f, P, [a,b]) \end{equation*}

for every partition \(P\) of \([a,b]\).

Suppose first that the \(\varepsilon\)-condition holds. Let \(\varepsilon > 0\) and choose a partition \(P\) of \([a,b]\) with \(U(f,P,[a,b]) - L(f,P,[a,b]) < \varepsilon\). Then, using the two inequalities above and 1.8,

\begin{equation*} 0 \le U(f,[a,b]) - L(f,[a,b]) \le U(f,P,[a,b]) - L(f,P,[a,b]) < \varepsilon . \end{equation*}

Since this holds for every \(\varepsilon > 0\), we get \(U(f,[a,b]) - L(f,[a,b]) = 0\). Thus \(L(f,[a,b]) = U(f,[a,b])\), which by 1.9 says exactly that \(f\) is Riemann integrable.

Conversely, suppose \(f\) is Riemann integrable, so \(L(f,[a,b]) = U(f,[a,b])\), and let \(\varepsilon > 0\). Because \(L(f,[a,b])\) is the supremum of the lower Riemann sums, the number \(L(f,[a,b]) - \frac{\varepsilon}{2}\) is not an upper bound for them, so there is a partition \(P_1\) of \([a,b]\) with

\begin{equation*} L(f, P_1, [a,b]) > L(f, [a,b]) - \tfrac{\varepsilon}{2} . \end{equation*}

Because \(U(f,[a,b])\) is the infimum of the upper Riemann sums, there is likewise a partition \(P_2\) of \([a,b]\) with

\begin{equation*} U(f, P_2, [a,b]) < U(f, [a,b]) + \tfrac{\varepsilon}{2} . \end{equation*}

Let \(P\) be the partition of \([a,b]\) obtained by merging the lists defining \(P_1\) and \(P_2\); then the lists defining \(P_1\) and \(P_2\) are both sublists of the list defining \(P\), so 1.5 gives

\begin{equation*} L(f, P, [a,b]) \ge L(f, P_1, [a,b]) \quad\text{and}\quad U(f, P, [a,b]) \le U(f, P_2, [a,b]). \end{equation*}

Combining these with the previous two displays and with \(L(f,[a,b]) = U(f,[a,b])\),

\begin{equation*} \begin{aligned} U(f, P, [a,b]) - L(f, P, [a,b]) &\le U(f, P_2, [a,b]) - L(f, P_1, [a,b]) \\ &< \Big( U(f,[a,b]) + \tfrac{\varepsilon}{2} \Big) - \Big( L(f,[a,b]) - \tfrac{\varepsilon}{2} \Big) \\ &= \varepsilon . \end{aligned} \end{equation*}

Thus \(P\) is a partition of the required kind, completing the proof.

Problem (1A.4)

Suppose \(f, g : [a, b] \to \mathbb{R}\) are Riemann integrable. Prove that \(f + g\) is Riemann integrable on \([a, b]\) and

\begin{equation*} \int_a^b (f + g) = \int_a^b f + \int_a^b g . \end{equation*}

Solution

Both \(f\) and \(g\) are bounded, hence so is \(f+g\); thus the Riemann sums of \(f+g\) are defined.

Step 1: a pointwise comparison of Riemann sums. Let \(A\) be a nonempty subset of \([a,b]\). For every \(x \in A\) we have \(f(x) + g(x) \ge \inf_A f + \inf_A g\), so \(\inf_A f + \inf_A g\) is a lower bound for \(\{(f+g)(x) : x \in A\}\) and therefore

\begin{equation*} \inf_A f + \inf_A g \le \inf_A (f+g). \end{equation*}

Symmetrically, \(\sup_A f + \sup_A g\) is an upper bound for that set, so

\begin{equation*} \sup_A (f+g) \le \sup_A f + \sup_A g . \end{equation*}

Applying these with \(A = [x_{j-1}, x_j]\) for each subinterval of an arbitrary partition \(P\) of \([a,b]\), multiplying by \(x_j - x_{j-1} > 0\), and summing over \(j\) gives

\begin{equation*} L(f, P, [a,b]) + L(g, P, [a,b]) \le L(f+g, P, [a,b]) \le U(f+g, P, [a,b]) \le U(f, P, [a,b]) + U(g, P, [a,b]), \end{equation*}

where the middle inequality holds because an infimum never exceeds the corresponding supremum.

Step 2: integrability. Let \(\varepsilon > 0\). Since \(f\) and \(g\) are Riemann integrable, Exercise 3 in this section gives partitions \(P_1\) and \(P_2\) of \([a,b]\) with

\begin{equation*} U(f, P_1, [a,b]) - L(f, P_1, [a,b]) < \tfrac{\varepsilon}{2}, \qquad U(g, P_2, [a,b]) - L(g, P_2, [a,b]) < \tfrac{\varepsilon}{2}. \end{equation*}

Let \(P\) be the partition obtained by merging the lists defining \(P_1\) and \(P_2\). By 1.5, passing from \(P_1\) to \(P\) can only increase the lower sum of \(f\) and decrease its upper sum, so \(U(f,P,[a,b]) - L(f,P,[a,b]) < \frac{\varepsilon}{2}\); for the same reason \(U(g,P,[a,b]) - L(g,P,[a,b]) < \frac{\varepsilon}{2}\). Now Step 1 gives

\begin{equation*} \begin{aligned} U(f+g, P, [a,b]) - L(f+g, P, [a,b]) &\le \big( U(f,P,[a,b]) + U(g,P,[a,b]) \big) - \big( L(f,P,[a,b]) + L(g,P,[a,b]) \big) \\ &< \tfrac{\varepsilon}{2} + \tfrac{\varepsilon}{2} = \varepsilon . \end{aligned} \end{equation*}

Since \(\varepsilon > 0\) was arbitrary, Exercise 3 shows that \(f + g\) is Riemann integrable on \([a,b]\).

Step 3: the value of the integral. Keep \(\varepsilon > 0\) and the partition \(P\) from Step 2. Using \(\int_a^b (f+g) = U(f+g,[a,b]) \le U(f+g,P,[a,b])\) together with Step 1,

\begin{equation*} \int_a^b (f+g) \le U(f,P,[a,b]) + U(g,P,[a,b]) < \Big( L(f,P,[a,b]) + \tfrac{\varepsilon}{2} \Big) + \Big( L(g,P,[a,b]) + \tfrac{\varepsilon}{2} \Big) \le \int_a^b f + \int_a^b g + \varepsilon, \end{equation*}

where the last step uses \(L(f,P,[a,b]) \le L(f,[a,b]) = \int_a^b f\) and similarly for \(g\). In the same way, using \(\int_a^b (f+g) = L(f+g,[a,b]) \ge L(f+g,P,[a,b])\) and Step 1,

\begin{equation*} \int_a^b (f+g) \ge L(f,P,[a,b]) + L(g,P,[a,b]) > \Big( U(f,P,[a,b]) - \tfrac{\varepsilon}{2} \Big) + \Big( U(g,P,[a,b]) - \tfrac{\varepsilon}{2} \Big) \ge \int_a^b f + \int_a^b g - \varepsilon , \end{equation*}

where the last step uses \(U(f,P,[a,b]) \ge U(f,[a,b]) = \int_a^b f\) and similarly for \(g\). Hence

\begin{equation*} \left| \int_a^b (f+g) - \Big( \int_a^b f + \int_a^b g \Big) \right| < \varepsilon \end{equation*}

for every \(\varepsilon > 0\), which forces \(\int_a^b (f+g) = \int_a^b f + \int_a^b g\).

Problem (1A.5)

Suppose \(f : [a, b] \to \mathbb{R}\) is Riemann integrable. Prove that the function \(-f\) is Riemann integrable on \([a, b]\) and

\begin{equation*} \int_a^b (-f) = - \int_a^b f . \end{equation*}

Solution

Since \(f\) is bounded, so is \(-f\).

Step 1: negation swaps infima and suprema. Let \(A\) be a nonempty subset of \([a,b]\) and put \(\alpha = \sup_A f\). For every \(x \in A\) we have \(f(x) \le \alpha\), hence \(-f(x) \ge -\alpha\), so \(-\alpha\) is a lower bound for \(\{-f(x) : x \in A\}\). If \(\beta\) is any lower bound for that set, then \(-\beta\) is an upper bound for \(\{f(x) : x \in A\}\), so \(\alpha \le -\beta\), i.e. \(\beta \le -\alpha\). Thus \(-\alpha\) is the greatest lower bound:

\begin{equation*} \inf_A (-f) = - \sup_A f . \end{equation*}

Replacing \(f\) by \(-f\) in this identity (and using \(-(-f) = f\)) gives \(\inf_A f = -\sup_A(-f)\), that is,

\begin{equation*} \sup_A (-f) = - \inf_A f . \end{equation*}

Step 2: negation swaps lower and upper Riemann sums. Let \(P\) be the partition \(a = x_0 < \cdots < x_n = b\). Applying Step 1 with \(A = [x_{j-1}, x_j]\),

\begin{equation*} L(-f, P, [a,b]) = \sum_{j=1}^n (x_j - x_{j-1}) \inf_{[x_{j-1},x_j]} (-f) = - \sum_{j=1}^n (x_j - x_{j-1}) \sup_{[x_{j-1},x_j]} f = -\,U(f, P, [a,b]), \end{equation*}

and symmetrically \(U(-f, P, [a,b]) = -\,L(f, P, [a,b])\).

Step 3: negation swaps lower and upper Riemann integrals. For any nonempty set \(S \subseteq \mathbb{R}\) we have \(\sup(-S) = -\inf S\) and \(\inf(-S) = -\sup S\), by the argument of Step 1 applied to \(S\). Taking \(S\) to be the set of upper Riemann sums of \(f\) over all partitions of \([a,b]\), Step 2 shows that the set of lower Riemann sums of \(-f\) is exactly \(-S\), so

\begin{equation*} L(-f, [a,b]) = \sup_P L(-f, P, [a,b]) = \sup(-S) = -\inf S = -\,U(f, [a,b]) . \end{equation*}

Likewise, taking \(T\) to be the set of lower Riemann sums of \(f\),

\begin{equation*} U(-f, [a,b]) = \inf_P U(-f, P, [a,b]) = \inf(-T) = -\sup T = -\,L(f, [a,b]) . \end{equation*}

Step 4: conclusion. Because \(f\) is Riemann integrable, \(L(f,[a,b]) = U(f,[a,b]) = \int_a^b f\). Hence Step 3 gives

\begin{equation*} L(-f, [a,b]) = -\int_a^b f = U(-f, [a,b]) . \end{equation*}

So the lower and upper Riemann integrals of \(-f\) agree; by 1.9, \(-f\) is Riemann integrable on \([a,b]\) and \(\int_a^b (-f) = -\int_a^b f\).

Problem (1A.6)

Suppose \(f : [a, b] \to \mathbb{R}\) is Riemann integrable. Suppose \(g : [a, b] \to \mathbb{R}\) is a function such that \(g(x) = f(x)\) for all except finitely many \(x \in [a, b]\). Prove that \(g\) is Riemann integrable on \([a, b]\) and

\begin{equation*} \int_a^b g = \int_a^b f . \end{equation*}

Solution

Define \(h : [a,b] \to \mathbb{R}\) by \(h = g - f\). Then \(h(x) = 0\) for all except finitely many \(x \in [a,b]\), and \(g = f + h\). By Exercise 4 in this section it suffices to prove that \(h\) is Riemann integrable on \([a,b]\) with \(\int_a^b h = 0\): granting that, \(g = f + h\) is Riemann integrable and \(\int_a^b g = \int_a^b f + \int_a^b h = \int_a^b f\).

So let \(c_1 < c_2 < \cdots < c_m\) be the points of \([a,b]\) at which \(h\) is nonzero (if there are none then \(h \equiv 0\), whose lower and upper Riemann sums are all \(0\), so \(h\) is Riemann integrable with integral \(0\) and we are done). Set

\begin{equation*} M = \max\{ |h(c_1)|, \ldots, |h(c_m)| \} > 0 . \end{equation*}

Every value of \(h\) lies in \([-M, M]\), so \(h\) is bounded.

Construction of good partitions. Let \(\delta > 0\). Let \(P_\delta\) be the partition of \([a,b]\) whose points are the distinct elements of

\begin{equation*} \{a, b\} \cup \big( \{ c_i - \delta,\; c_i + \delta \;:\; i = 1, \ldots, m \} \cap (a,b) \big), \end{equation*}

listed in increasing order. This is a genuine partition of \([a,b]\): it is a finite increasing list of points of \([a,b]\) beginning at \(a\) and ending at \(b\).

Claim: if \([u,w]\) is a subinterval of \(P_\delta\) and \(c_i \in [u,w]\), then \([u,w] \subseteq [c_i - \delta, c_i + \delta]\). Indeed, suppose \(u < c_i - \delta\). Since also \(c_i - \delta < c_i \le w\), the number \(c_i - \delta\) lies in \((u,w)\), and since \(a \le u < c_i - \delta < w \le b\) it lies in \((a,b)\). Hence \(c_i - \delta\) is one of the points of \(P_\delta\) strictly between the consecutive partition points \(u\) and \(w\), which is impossible. So \(u \ge c_i - \delta\), and the symmetric argument gives \(w \le c_i + \delta\), proving the claim.

Estimating the Riemann sums. Call a subinterval of \(P_\delta\) bad if it contains at least one of \(c_1, \ldots, c_m\), and good otherwise. On a good subinterval \(I\) we have \(h \equiv 0\), so \(\sup_I h = \inf_I h = 0\) and that subinterval contributes \(0\) to both \(L(h, P_\delta, [a,b])\) and \(U(h, P_\delta, [a,b])\). On a bad subinterval \(I\) the values of \(h\) lie in \([-M,M]\), so \(|\sup_I h| \le M\) and \(|\inf_I h| \le M\).

For a fixed \(i\), the bad subintervals containing \(c_i\) (there are at most two of them, since \(c_i\) can be an endpoint shared by two subintervals) are all contained in \([c_i - \delta, c_i + \delta]\) by the claim, and distinct subintervals of a partition have disjoint interiors; hence their lengths sum to at most \(2\delta\). Since every bad subinterval contains some \(c_i\), the total length of the bad subintervals is at most \(2m\delta\). Therefore

\begin{equation*} |U(h, P_\delta, [a,b])| \le 2m\delta M \quad\text{and}\quad |L(h, P_\delta, [a,b])| \le 2m\delta M . \end{equation*}

Conclusion. Let \(\varepsilon > 0\) and take \(\delta = \dfrac{\varepsilon}{2mM}\) above. Then

\begin{equation*} U(h, [a,b]) \le U(h, P_\delta, [a,b]) \le \varepsilon \quad\text{and}\quad L(h, [a,b]) \ge L(h, P_\delta, [a,b]) \ge -\varepsilon . \end{equation*}

As \(\varepsilon > 0\) was arbitrary, \(U(h,[a,b]) \le 0\) and \(L(h,[a,b]) \ge 0\). Combined with \(L(h,[a,b]) \le U(h,[a,b])\) from 1.8, this gives

\begin{equation*} L(h, [a,b]) = U(h, [a,b]) = 0 . \end{equation*}

Thus \(h\) is Riemann integrable on \([a,b]\) with \(\int_a^b h = 0\), and by the first paragraph \(g\) is Riemann integrable on \([a,b]\) with \(\int_a^b g = \int_a^b f\).

Problem (1A.7)

Suppose \(f : [a, b] \to \mathbb{R}\) is a bounded function. For \(n \in \mathbb{Z}^{+}\), let \(P_n\) denote the partition that divides \([a, b]\) into \(2^n\) intervals of equal size. Prove that

\begin{equation*} L(f, [a, b]) = \lim_{n \to \infty} L(f, P_n, [a, b]) \quad \text{and} \quad U(f, [a, b]) = \lim_{n \to \infty} U(f, P_n, [a, b]). \end{equation*}

Solution

Write \(P_n\) for the partition with points \(x_j^{(n)} = a + \dfrac{j(b-a)}{2^n}\) for \(j = 0, 1, \ldots, 2^n\); each of its \(2^n\) subintervals has length \(\dfrac{b-a}{2^n}\). Since \(f\) is bounded we may set \(M = \sup_{[a,b]} |f| < \infty\), so every value of \(f\) lies in \([-M, M]\).

A refinement estimate. Suppose \(Q\) is a partition of \([a,b]\) all of whose subintervals have length at most \(\ell\), and \(Q^{\prime}\) is a partition whose defining list is obtained from that of \(Q\) by adjoining \(p\) additional points. Then

\begin{equation*} L(f, Q^{\prime}, [a,b]) \le L(f, Q, [a,b]) + 2pM\ell . \end{equation*}

To see this, first take \(p = 1\), say the adjoined point is \(v\). If \(v\) is already a point of \(Q\) the two sums are equal and there is nothing to prove. Otherwise \(v\) lies in the interior of exactly one subinterval \([u,w]\) of \(Q\), and the two lower sums have identical terms except for those coming from \([u,w]\), so

\begin{equation*} \begin{aligned} L(f, Q^{\prime}, [a,b]) - L(f, Q, [a,b]) &= (v-u) \inf_{[u,v]} f + (w-v) \inf_{[v,w]} f - (w-u) \inf_{[u,w]} f \\ &\le (v-u)M + (w-v)M + (w-u)M \\ &= 2M(w-u) \le 2M\ell , \end{aligned} \end{equation*}

using \(|\inf_A f| \le M\) for every nonempty \(A \subseteq [a,b]\). For general \(p\), adjoin the \(p\) points one at a time; every intermediate partition is a refinement of \(Q\), so all of its subintervals still have length at most \(\ell\), and applying the case \(p = 1\) at each of the \(p\) steps and summing gives the estimate.

The lower integral. Since \(x_j^{(n)} = a + \dfrac{2j(b-a)}{2^{n+1}} = x_{2j}^{(n+1)}\), the list defining \(P_n\) is a sublist of the list defining \(P_{n+1}\). Hence 1.5 gives

\begin{equation*} L(f, P_n, [a,b]) \le L(f, P_{n+1}, [a,b]) \quad \text{for all } n \in \mathbb{Z}^{+}, \end{equation*}

and by the definition 1.7 of the lower Riemann integral,

\begin{equation*} L(f, P_n, [a,b]) \le L(f, [a,b]) \quad \text{for all } n \in \mathbb{Z}^{+}. \end{equation*}

So \(\big( L(f, P_n, [a,b]) \big)_{n \in \mathbb{Z}^{+}}\) is an increasing sequence bounded above by \(L(f,[a,b])\); therefore it converges and

\begin{equation*} \lim_{n \to \infty} L(f, P_n, [a,b]) \le L(f, [a,b]) . \end{equation*}

For the reverse inequality, let \(\varepsilon > 0\). Because \(L(f,[a,b])\) is the supremum of the lower Riemann sums, there is a partition \(P\) of \([a,b]\) with

\begin{equation*} L(f, P, [a,b]) > L(f, [a,b]) - \varepsilon . \end{equation*}

Say \(P\) is the list \(a = y_0 < y_1 < \cdots < y_k = b\), so \(P\) has \(k - 1\) points other than \(a\) and \(b\). For each \(n\) let \(Q_n\) be the partition obtained by merging the lists defining \(P_n\) and \(P\). Then the list defining \(Q_n\) is obtained from that of \(P_n\) by adjoining at most \(k-1\) points, and every subinterval of \(P_n\) has length \(\dfrac{b-a}{2^n}\); so the refinement estimate gives

\begin{equation*} L(f, Q_n, [a,b]) \le L(f, P_n, [a,b]) + \frac{2(k-1)M(b-a)}{2^n} . \end{equation*}

On the other hand, the list defining \(P\) is a sublist of the list defining \(Q_n\), so \(L(f, P, [a,b]) \le L(f, Q_n, [a,b])\) by 1.5. Combining the last three displays,

\begin{equation*} L(f, P_n, [a,b]) \ge L(f, P, [a,b]) - \frac{2(k-1)M(b-a)}{2^n} > L(f, [a,b]) - \varepsilon - \frac{2(k-1)M(b-a)}{2^n} . \end{equation*}

Choose \(N \in \mathbb{Z}^{+}\) with \(\dfrac{2(k-1)M(b-a)}{2^N} < \varepsilon\) (possible because \(2^n \to \infty\)). Then for all \(n \ge N\),

\begin{equation*} L(f, [a,b]) - 2\varepsilon < L(f, P_n, [a,b]) \le L(f, [a,b]) . \end{equation*}

Hence \(\big| \lim_{n \to \infty} L(f, P_n, [a,b]) - L(f,[a,b]) \big| \le 2\varepsilon\) for every \(\varepsilon > 0\), and therefore

\begin{equation*} \lim_{n \to \infty} L(f, P_n, [a,b]) = L(f, [a,b]) . \end{equation*}

The upper integral. Apply the result just proved to the bounded function \(-f\), and use the identities established in Steps 2 and 3 of Exercise 5 in this section, whose proofs used only that the function is bounded:

\begin{equation*} L(-f, P, [a,b]) = -\,U(f, P, [a,b]) \ \text{ for every partition } P, \qquad L(-f, [a,b]) = -\,U(f, [a,b]) . \end{equation*}

Therefore

\begin{equation*} U(f, [a,b]) = -\,L(-f, [a,b]) = -\lim_{n \to \infty} L(-f, P_n, [a,b]) = \lim_{n \to \infty} \big( -L(-f, P_n, [a,b]) \big) = \lim_{n \to \infty} U(f, P_n, [a,b]), \end{equation*}

as desired.

Problem (1A.8)

Suppose \(f \colon [a,b] \to \mathbb{R}\) is Riemann integrable. Prove that

\begin{equation*} \int_a^b f = \lim_{n \to \infty} \frac{b-a}{n} \sum_{j=1}^{n} f\!\left(a + \tfrac{j(b-a)}{n}\right). \end{equation*}

Solution

For \(n \in \mathbb{Z}^+\) let \(P_n\) denote the equally spaced partition \(x_0, x_1, \dots, x_n\) of \([a,b]\), where

\begin{equation*} x_j = a + \frac{j(b-a)}{n} \qquad \text{for } j = 0, 1, \dots, n, \end{equation*}

so that \(x_j - x_{j-1} = \frac{b-a}{n}\) for each \(j\). Write

\begin{equation*} R_n = \frac{b-a}{n} \sum_{j=1}^{n} f(x_j), \end{equation*}

which is the quantity whose limit we must compute.

Step 1: \(R_n\) is trapped between the lower and upper Riemann sums for \(P_n\).

For each \(j \in \{1,\dots,n\}\) we have \(x_j \in [x_{j-1},x_j]\), and therefore

\begin{equation*} \inf_{[x_{j-1},x_j]} f \;\le\; f(x_j) \;\le\; \sup_{[x_{j-1},x_j]} f . \end{equation*}

Multiplying by \(x_j - x_{j-1} = \frac{b-a}{n} > 0\) and summing over \(j\) gives

\begin{equation*} L(f, P_n, [a,b]) \;\le\; R_n \;\le\; U(f, P_n, [a,b]). \end{equation*}

So it suffices to prove that both \(L(f,P_n,[a,b])\) and \(U(f,P_n,[a,b])\) converge to \(\int_a^b f\) as \(n \to \infty\); the squeeze then finishes the proof.

Step 2: \(U(f,P_n,[a,b]) \to U(f,[a,b])\).

Because \(f\) is Riemann integrable it is by definition bounded, so we may set

\begin{equation*} M = \sup_{[a,b]} |f| < \infty . \end{equation*}

Let \(\varepsilon > 0\). By the definition 1.7 of the upper Riemann integral as an infimum over partitions, there is a partition \(P\) of \([a,b]\) with

\begin{equation*} U(f, P, [a,b]) < U(f,[a,b]) + \varepsilon . \end{equation*}

Say \(P\) is the list \(y_0, y_1, \dots, y_m\); thus \(P\) has at most \(m-1\) points lying in the open interval \((a,b)\).

Fix \(n \in \mathbb{Z}^+\) and let \(P_n^{\prime}\) be the partition obtained by merging the lists defining \(P_n\) and \(P\). Since the list defining \(P_n\) is a sublist of the list defining \(P_n^{\prime}\), and likewise for \(P\), two applications of 1.5 give

\begin{equation*} U(f, P_n^{\prime}, [a,b]) \le U(f, P, [a,b]) \qquad \text{and} \qquad L(f,P,[a,b]) \le L(f,P_n^{\prime},[a,b]). \end{equation*}

We now compare \(U(f,P_n,[a,b])\) with \(U(f,P_n^{\prime},[a,b])\). The partition \(P_n^{\prime}\) is obtained from \(P_n\) by adjoining those points of \(P\) that are not already in \(P_n\); there are at most \(m-1\) such points (only the points \(y_1,\dots,y_{m-1}\) of \((a,b)\) can be new, since \(y_0 = a\) and \(y_m = b\) already belong to \(P_n\)). Call a subinterval \([x_{j-1},x_j]\) of \(P_n\) affected if at least one new point lies in its interior; at most \(m-1\) subintervals are affected.

If \([x_{j-1},x_j]\) is not affected, then \([x_{j-1},x_j]\) is still a subinterval of \(P_n^{\prime}\) and it contributes the identical term \((x_j - x_{j-1}) \sup_{[x_{j-1},x_j]} f\) to both upper sums. If \([x_{j-1},x_j]\) is affected, it is subdivided in \(P_n^{\prime}\) into consecutive subintervals \([t_{i-1},t_i]\), \(i = 1,\dots,r\), with \(t_0 = x_{j-1}\) and \(t_r = x_j\). Its contribution to \(U(f,P_n,[a,b])\) is at most \((x_j - x_{j-1}) M\), while its total contribution to \(U(f,P_n^{\prime},[a,b])\) is

\begin{equation*} \sum_{i=1}^{r} (t_i - t_{i-1}) \sup_{[t_{i-1},t_i]} f \;\ge\; \sum_{i=1}^{r} (t_i - t_{i-1})(-M) \;=\; -(x_j - x_{j-1}) M . \end{equation*}

Hence the affected subinterval \([x_{j-1},x_j]\) contributes at most \(2M(x_j - x_{j-1}) = \frac{2M(b-a)}{n}\) more to \(U(f,P_n,[a,b])\) than to \(U(f,P_n^{\prime},[a,b])\). Summing over the at most \(m-1\) affected subintervals,

\begin{equation*} U(f,P_n,[a,b]) \;\le\; U(f,P_n^{\prime},[a,b]) + \frac{2M(m-1)(b-a)}{n} \;\le\; U(f,P,[a,b]) + \frac{2M(m-1)(b-a)}{n}. \end{equation*}

Combining with the choice of \(P\),

\begin{equation*} U(f,P_n,[a,b]) \;<\; U(f,[a,b]) + \varepsilon + \frac{2M(m-1)(b-a)}{n}. \end{equation*}

Note that \(m\) and \(M\) do not depend on \(n\). Thus there exists \(N \in \mathbb{Z}^+\) such that \(\frac{2M(m-1)(b-a)}{n} < \varepsilon\) for all \(n \ge N\), and then

\begin{equation*} U(f,[a,b]) \;\le\; U(f,P_n,[a,b]) \;<\; U(f,[a,b]) + 2\varepsilon \qquad \text{for all } n \ge N, \end{equation*}

where the first inequality holds because \(U(f,[a,b])\) is the infimum of the upper sums (1.7). As \(\varepsilon > 0\) was arbitrary, this proves

\begin{equation*} \lim_{n \to \infty} U(f,P_n,[a,b]) = U(f,[a,b]). \end{equation*}

Step 3: \(L(f,P_n,[a,b]) \to L(f,[a,b])\).

The argument is symmetric. Choosing \(P\) with \(L(f,P,[a,b]) > L(f,[a,b]) - \varepsilon\) and estimating in the same way (each affected subinterval contributes at most \(\frac{2M(b-a)}{n}\) less to \(L(f,P_n,[a,b])\) than to \(L(f,P_n^{\prime},[a,b])\), because its contribution to \(L(f,P_n,[a,b])\) is at least \(-(x_j - x_{j-1})M\) and its contribution to \(L(f,P_n^{\prime},[a,b])\) is at most \((x_j-x_{j-1})M\)) yields

\begin{equation*} L(f,P_n,[a,b]) \;\ge\; L(f,P_n^{\prime},[a,b]) - \frac{2M(m-1)(b-a)}{n} \;\ge\; L(f,P,[a,b]) - \frac{2M(m-1)(b-a)}{n}, \end{equation*}

and hence \(L(f,[a,b]) - 2\varepsilon < L(f,P_n,[a,b]) \le L(f,[a,b])\) for all large \(n\). Thus \(L(f,P_n,[a,b]) \to L(f,[a,b])\).

Step 4: Conclusion.

Because \(f\) is Riemann integrable, 1.9 gives

\begin{equation*} L(f,[a,b]) = U(f,[a,b]) = \int_a^b f . \end{equation*}

By Steps 2 and 3 both \(L(f,P_n,[a,b])\) and \(U(f,P_n,[a,b])\) converge to \(\int_a^b f\). Since Step 1 gives

\begin{equation*} L(f,P_n,[a,b]) \le R_n \le U(f,P_n,[a,b]), \end{equation*}

the squeeze theorem yields \(R_n \to \int_a^b f\), which is exactly the assertion

\begin{equation*} \int_a^b f = \lim_{n \to \infty} \frac{b-a}{n} \sum_{j=1}^{n} f\!\left(a + \tfrac{j(b-a)}{n}\right). \qquad \blacksquare \end{equation*}

Problem (1A.9)

Suppose \(f \colon [a,b] \to \mathbb{R}\) is Riemann integrable. Prove that if \(c, d \in \mathbb{R}\) and \(a \le c < d \le b\), then \(f\) is Riemann integrable on \([c,d]\).

[To say that \(f\) is Riemann integrable on \([c,d]\) means that \(f\) with its domain restricted to \([c,d]\) is Riemann integrable.]

Solution

Throughout, we use the following criterion, which is Exercise 3 in Section 1A. Because it is used repeatedly below and in later exercises, we prove it here.

Riemann criterion. A bounded function \(g \colon [\alpha,\beta] \to \mathbb{R}\) is Riemann integrable if and only if for every \(\varepsilon > 0\) there is a partition \(P\) of \([\alpha,\beta]\) with

\begin{equation*} U(g,P,[\alpha,\beta]) - L(g,P,[\alpha,\beta]) < \varepsilon . \end{equation*}

Proof of the criterion. Suppose first that such partitions exist. Fix \(\varepsilon > 0\) and choose \(P\) as above. By 1.7 we have \(U(g,[\alpha,\beta]) \le U(g,P,[\alpha,\beta])\) and \(L(g,[\alpha,\beta]) \ge L(g,P,[\alpha,\beta])\), so

\begin{equation*} 0 \le U(g,[\alpha,\beta]) - L(g,[\alpha,\beta]) \le U(g,P,[\alpha,\beta]) - L(g,P,[\alpha,\beta]) < \varepsilon, \end{equation*}

the first inequality coming from 1.8. Since \(\varepsilon > 0\) is arbitrary, \(L(g,[\alpha,\beta]) = U(g,[\alpha,\beta])\), so \(g\) is Riemann integrable by 1.9.

Conversely, suppose \(g\) is Riemann integrable and let \(\varepsilon > 0\). By 1.7 there are partitions \(P_1, P_2\) of \([\alpha,\beta]\) with

\begin{equation*} U(g,P_1,[\alpha,\beta]) < U(g,[\alpha,\beta]) + \tfrac{\varepsilon}{2}, \qquad L(g,P_2,[\alpha,\beta]) > L(g,[\alpha,\beta]) - \tfrac{\varepsilon}{2}. \end{equation*}

Let \(P\) be the partition obtained by merging the lists defining \(P_1\) and \(P_2\). By 1.5,

\begin{equation*} U(g,P,[\alpha,\beta]) \le U(g,P_1,[\alpha,\beta]) \quad \text{and} \quad L(g,P,[\alpha,\beta]) \ge L(g,P_2,[\alpha,\beta]). \end{equation*}

Since \(U(g,[\alpha,\beta]) = L(g,[\alpha,\beta])\), subtracting gives

\begin{equation*} U(g,P,[\alpha,\beta]) - L(g,P,[\alpha,\beta]) < \big(U(g,[\alpha,\beta]) + \tfrac{\varepsilon}{2}\big) - \big(L(g,[\alpha,\beta]) - \tfrac{\varepsilon}{2}\big) = \varepsilon, \end{equation*}

which completes the proof of the criterion.

Now we prove the exercise. Since \(f\) is bounded on \([a,b]\), its restriction to \([c,d] \subseteq [a,b]\) is bounded, so the lower and upper Riemann sums and integrals of \(f\) on \([c,d]\) are defined.

Let \(\varepsilon > 0\). Since \(f\) is Riemann integrable on \([a,b]\), the criterion supplies a partition \(P\) of \([a,b]\) with

\begin{equation*} U(f,P,[a,b]) - L(f,P,[a,b]) < \varepsilon . \end{equation*}

Let \(P^{\prime}\) be the partition of \([a,b]\) obtained by adjoining the points \(c\) and \(d\) to the list defining \(P\) (they may already be present). The list defining \(P\) is a sublist of the list defining \(P^{\prime}\), so 1.5 gives

\begin{equation*} U(f,P^{\prime},[a,b]) \le U(f,P,[a,b]) \quad \text{and} \quad L(f,P^{\prime},[a,b]) \ge L(f,P,[a,b]), \end{equation*}

and therefore

\begin{equation*} U(f,P^{\prime},[a,b]) - L(f,P^{\prime},[a,b]) \le U(f,P,[a,b]) - L(f,P,[a,b]) < \varepsilon . \end{equation*}

Write \(P^{\prime}\) as the list \(x_0, x_1, \dots, x_n\) with \(a = x_0 < x_1 < \cdots < x_n = b\). Since \(c\) and \(d\) appear in this list and \(c < d\), there are indices \(p < q\) with \(x_p = c\) and \(x_q = d\). Then

\begin{equation*} Q : \quad x_p, x_{p+1}, \dots, x_q \end{equation*}

is a partition of \([c,d]\).

For each \(j \in \{p+1, \dots, q\}\) the subinterval \([x_{j-1},x_j]\) is contained in \([c,d]\), so the supremum and infimum of \(f\) over \([x_{j-1},x_j]\) are unchanged when \(f\) is restricted to \([c,d]\). Hence, writing \(f\) also for its restriction to \([c,d]\),

\begin{equation*} \begin{aligned} U(f,Q,[c,d]) - L(f,Q,[c,d]) &= \sum_{j=p+1}^{q} (x_j - x_{j-1})\Big( \sup_{[x_{j-1},x_j]} f - \inf_{[x_{j-1},x_j]} f \Big) \\ &\le \sum_{j=1}^{n} (x_j - x_{j-1})\Big( \sup_{[x_{j-1},x_j]} f - \inf_{[x_{j-1},x_j]} f \Big) \\ &= U(f,P^{\prime},[a,b]) - L(f,P^{\prime},[a,b]) \\ &< \varepsilon , \end{aligned} \end{equation*}

where the inequality in the second line holds because every omitted term is nonnegative (each factor \(x_j - x_{j-1}\) is positive and each supremum is at least the corresponding infimum).

Thus for every \(\varepsilon > 0\) there is a partition \(Q\) of \([c,d]\) with \(U(f,Q,[c,d]) - L(f,Q,[c,d]) < \varepsilon\). By the criterion, \(f\) is Riemann integrable on \([c,d]\). \(\blacksquare\)

Problem (1A.10)

Suppose \(f \colon [a,b] \to \mathbb{R}\) is a bounded function and \(c \in (a,b)\). Prove that \(f\) is Riemann integrable on \([a,b]\) if and only if \(f\) is Riemann integrable on \([a,c]\) and \(f\) is Riemann integrable on \([c,b]\). Furthermore, prove that if these conditions hold, then

\begin{equation*} \int_a^b f = \int_a^c f + \int_c^b f . \end{equation*}

Solution

Since \(f\) is bounded on \([a,b]\), its restrictions to \([a,c]\) and to \([c,b]\) are bounded, so all four lower and upper Riemann integrals appearing below are defined. We write \(f\) for each of these restrictions as well.

Splitting a partition. Suppose \(P\) is a partition \(x_0, x_1, \dots, x_n\) of \([a,b]\) whose list contains \(c\), say \(c = x_p\) with \(0 < p < n\). Then

\begin{equation*} P_1 : \; x_0, \dots, x_p \quad \text{and} \quad P_2 : \; x_p, \dots, x_n \end{equation*}

are partitions of \([a,c]\) and \([c,b]\) respectively. Each subinterval \([x_{j-1},x_j]\) with \(j \le p\) is contained in \([a,c]\) and each with \(j > p\) is contained in \([c,b]\), so the relevant infima and suprema are unaffected by restricting the domain of \(f\). Splitting the defining sums in 1.3 at \(j = p\) therefore gives

\begin{equation*} L(f,P,[a,b]) = L(f,P_1,[a,c]) + L(f,P_2,[c,b]), \qquad U(f,P,[a,b]) = U(f,P_1,[a,c]) + U(f,P_2,[c,b]). \qquad (\ast) \end{equation*}

Conversely, if \(P_1\) is any partition of \([a,c]\) and \(P_2\) is any partition of \([c,b]\), then concatenating their lists (identifying the shared endpoint \(c\)) produces a partition \(P\) of \([a,b]\) containing \(c\), and \((\ast)\) again holds.

Step 1: additivity of the lower Riemann integral, namely

\begin{equation*} L(f,[a,b]) = L(f,[a,c]) + L(f,[c,b]). \end{equation*}

For the inequality \(\ge\): let \(P_1\) be a partition of \([a,c]\) and \(P_2\) a partition of \([c,b]\), and let \(P\) be their concatenation. By \((\ast)\) and 1.7,

\begin{equation*} L(f,P_1,[a,c]) + L(f,P_2,[c,b]) = L(f,P,[a,b]) \le L(f,[a,b]). \end{equation*}

Taking the supremum over \(P_1\) and then over \(P_2\) (the two are chosen independently) gives \(L(f,[a,c]) + L(f,[c,b]) \le L(f,[a,b])\).

For the inequality \(\le\): let \(P\) be any partition of \([a,b]\) and let \(P^{\prime}\) be obtained by adjoining \(c\) to the list defining \(P\). By 1.5 and then \((\ast)\) and 1.7,

\begin{equation*} L(f,P,[a,b]) \le L(f,P^{\prime},[a,b]) = L(f,P^{\prime}_1,[a,c]) + L(f,P^{\prime}_2,[c,b]) \le L(f,[a,c]) + L(f,[c,b]). \end{equation*}

Taking the supremum over all partitions \(P\) of \([a,b]\) gives \(L(f,[a,b]) \le L(f,[a,c]) + L(f,[c,b])\). This proves Step 1.

Step 2: additivity of the upper Riemann integral, namely

\begin{equation*} U(f,[a,b]) = U(f,[a,c]) + U(f,[c,b]). \end{equation*}

The proof is the mirror image of Step 1. For \(\le\): given partitions \(P_1\) of \([a,c]\) and \(P_2\) of \([c,b]\) with concatenation \(P\), we have \(U(f,[a,b]) \le U(f,P,[a,b]) = U(f,P_1,[a,c]) + U(f,P_2,[c,b])\); take the infimum over \(P_1\) and over \(P_2\). For \(\ge\): given a partition \(P\) of \([a,b]\), adjoin \(c\) to get \(P^{\prime}\) and use 1.5 to obtain

\begin{equation*} U(f,P,[a,b]) \ge U(f,P^{\prime},[a,b]) = U(f,P^{\prime}_1,[a,c]) + U(f,P^{\prime}_2,[c,b]) \ge U(f,[a,c]) + U(f,[c,b]); \end{equation*}

take the infimum over \(P\).

Step 3: the equivalence. Subtracting the identity of Step 1 from that of Step 2 gives

\begin{equation*} U(f,[a,b]) - L(f,[a,b]) = \Big( U(f,[a,c]) - L(f,[a,c]) \Big) + \Big( U(f,[c,b]) - L(f,[c,b]) \Big). \end{equation*}

By 1.8 each of the three differences above is nonnegative. Hence the left side equals \(0\) if and only if both terms on the right equal \(0\). By 1.9, the left side equals \(0\) exactly when \(f\) is Riemann integrable on \([a,b]\), the first term on the right equals \(0\) exactly when \(f\) is Riemann integrable on \([a,c]\), and the second exactly when \(f\) is Riemann integrable on \([c,b]\). This is precisely the asserted equivalence.

Step 4: the formula. Suppose the conditions hold. Then by 1.9 and Step 1,

\begin{equation*} \int_a^b f = L(f,[a,b]) = L(f,[a,c]) + L(f,[c,b]) = \int_a^c f + \int_c^b f . \qquad \blacksquare \end{equation*}

Problem (1A.11)

Suppose \(f \colon [a,b] \to \mathbb{R}\) is Riemann integrable. Define \(F \colon [a,b] \to \mathbb{R}\) by

\begin{equation*} F(t) = \begin{cases} 0 & \text{if } t = a, \\[2pt] \displaystyle\int_a^t f & \text{if } t \in (a,b]. \end{cases} \end{equation*}

Prove that \(F\) is continuous on \([a,b]\).

Solution

First note that \(F\) is well defined: if \(t \in (a,b]\) then \(a \le a < t \le b\), so by Exercise 9 in Section 1A the function \(f\) is Riemann integrable on \([a,t]\) and hence \(\int_a^t f\) makes sense.

Because \(f\) is Riemann integrable it is bounded, so we may set

\begin{equation*} M = \sup_{[a,b]} |f| < \infty . \end{equation*}

We prove the following Lipschitz estimate, which immediately gives continuity.

Claim. If \(s, t \in [a,b]\) with \(s < t\), then \(|F(t) - F(s)| \le M(t-s)\).

Proof of the claim. We first show that

\begin{equation*} F(t) - F(s) = \int_s^t f . \qquad (\dagger) \end{equation*}

If \(s = a\), then \(F(s) = 0\) and \(F(t) = \int_a^t f = \int_s^t f\), so \((\dagger)\) holds. If \(s > a\), then \(a < s < t\), so \(s \in (a,t)\). By Exercise 9 in Section 1A the function \(f\) is Riemann integrable on \([a,t]\) (and in particular bounded there), and by Exercise 10 in Section 1A applied on the interval \([a,t]\) with interior point \(s\),

\begin{equation*} \int_a^t f = \int_a^s f + \int_s^t f , \end{equation*}

which is exactly \(F(t) = F(s) + \int_s^t f\), proving \((\dagger)\).

Now \(f\) is Riemann integrable on \([s,t]\) by Exercise 9 in Section 1A, so 1.13 applies on \([s,t]\) and gives

\begin{equation*} (t-s) \inf_{[s,t]} f \;\le\; \int_s^t f \;\le\; (t-s) \sup_{[s,t]} f . \end{equation*}

Since \([s,t] \subseteq [a,b]\) we have \(-M \le \inf_{[s,t]} f\) and \(\sup_{[s,t]} f \le M\), hence

\begin{equation*} -M(t-s) \le \int_s^t f \le M(t-s), \qquad \text{that is,} \qquad \left| \int_s^t f \right| \le M(t-s). \end{equation*}

Combining with \((\dagger)\) proves the claim.

Note that the claim is symmetric in \(s\) and \(t\): it says that \(|F(t) - F(s)| \le M|t-s|\) for all \(s, t \in [a,b]\), since for \(s = t\) both sides are \(0\) and for \(s > t\) we may interchange the names of \(s\) and \(t\).

Continuity now follows. Let \(\varepsilon > 0\) and put \(\delta = \dfrac{\varepsilon}{M+1} > 0\). If \(s, t \in [a,b]\) satisfy \(|s-t| < \delta\), then

\begin{equation*} |F(t) - F(s)| \le M|t-s| \le M\delta = \frac{M}{M+1}\,\varepsilon < \varepsilon, \end{equation*}

where the final inequality holds because \(0 \le M < M+1\). (The nonstrict inequalities matter here: if \(M = 0\) the first two steps are equalities.) Thus \(F\) is uniformly continuous on \([a,b]\), and in particular continuous at every point of \([a,b]\). \(\blacksquare\)

Problem (1A.12)

Suppose \(f \colon [a,b] \to \mathbb{R}\) is Riemann integrable. Prove that \(|f|\) is Riemann integrable and that

\begin{equation*} \left| \int_a^b f \right| \le \int_a^b |f| . \end{equation*}

Solution

Step 1: an oscillation inequality. Let \(A\) be a nonempty subset of \([a,b]\). We claim

\begin{equation*} \sup_A |f| - \inf_A |f| \;\le\; \sup_A f - \inf_A f . \end{equation*}

Indeed, for any real-valued bounded function \(g\) on \(A\) we have

\begin{equation*} \sup_A g - \inf_A g = \sup\{ g(x) - g(y) : x, y \in A \}, \end{equation*}

because \(\sup_A g - \inf_A g = \sup_{x \in A} g(x) + \sup_{y \in A}(-g(y))\) and the supremum of a sum over independent variables is the sum of the suprema. Applying this with \(g = |f|\) and with \(g = f\), the claim reduces to showing that

\begin{equation*} |f|(x) - |f|(y) \le \sup\{ f(u) - f(v) : u,v \in A \} \qquad \text{for all } x,y \in A, \end{equation*}

which holds because the reverse triangle inequality gives

\begin{equation*} |f(x)| - |f(y)| \le \big| |f(x)| - |f(y)| \big| \le |f(x) - f(y)| = \max\{ f(x)-f(y), \, f(y)-f(x) \}, \end{equation*}

and both \(f(x)-f(y)\) and \(f(y)-f(x)\) are of the form \(f(u)-f(v)\) with \(u,v \in A\).

Step 2: \(|f|\) is Riemann integrable. Since \(f\) is bounded, so is \(|f|\), so the lower and upper Riemann sums of \(|f|\) are defined. Let \(P\) be an arbitrary partition \(x_0, \dots, x_n\) of \([a,b]\). Applying Step 1 with \(A = [x_{j-1},x_j]\) for each \(j\) (each such \(A\) is nonempty) and multiplying by \(x_j - x_{j-1} > 0\), then summing over \(j\), gives by 1.3

\begin{equation*} U(|f|,P,[a,b]) - L(|f|,P,[a,b]) \;\le\; U(f,P,[a,b]) - L(f,P,[a,b]). \end{equation*}

Let \(\varepsilon > 0\). By the Riemann criterion (Exercise 3 in Section 1A, proved in the solution to Exercise 9), there is a partition \(P\) of \([a,b]\) with \(U(f,P,[a,b]) - L(f,P,[a,b]) < \varepsilon\); the displayed inequality then gives \(U(|f|,P,[a,b]) - L(|f|,P,[a,b]) < \varepsilon\). By the same criterion, \(|f|\) is Riemann integrable.

Step 3: monotonicity of the Riemann integral. Suppose \(g, h \colon [a,b] \to \mathbb{R}\) are Riemann integrable with \(g(x) \le h(x)\) for all \(x \in [a,b]\). Then for every subinterval \([x_{j-1},x_j]\) we have \(\inf_{[x_{j-1},x_j]} g \le \inf_{[x_{j-1},x_j]} h\), so \(L(g,P,[a,b]) \le L(h,P,[a,b])\) for every partition \(P\) of \([a,b]\). Taking the supremum over \(P\) and using 1.7 and 1.9,

\begin{equation*} \int_a^b g = L(g,[a,b]) \le L(h,[a,b]) = \int_a^b h . \end{equation*}

Step 4: the inequality. By Step 2 the function \(|f|\) is Riemann integrable, and by Exercise 5 in Section 1A the function \(-|f|\) is Riemann integrable with

\begin{equation*} \int_a^b (-|f|) = - \int_a^b |f| . \end{equation*}

For every \(x \in [a,b]\) we have \(-|f(x)| \le f(x) \le |f(x)|\). Applying Step 3 to the pair \((f, |f|)\) and to the pair \((-|f|, f)\) gives

\begin{equation*}

  • \int_a^b |f| = \int_a^b (-|f|) \;\le\; \int_a^b f \;\le\; \int_a^b |f| . \end{equation*}

A real number lying between \(-r\) and \(r\) has absolute value at most \(r\), so

\begin{equation*} \left| \int_a^b f \right| \le \int_a^b |f| . \qquad \blacksquare \end{equation*}

Problem (1A.13)

Suppose \(f \colon [a,b] \to \mathbb{R}\) is an increasing function, meaning that \(c, d \in [a,b]\) with \(c < d\) implies \(f( c) \le f(d)\). Prove that \(f\) is Riemann integrable on \([a,b]\).

Solution

Step 1: \(f\) is bounded. For every \(x \in [a,b]\) we have \(a \le x \le b\), so monotonicity gives \(f(a) \le f(x) \le f(b)\). Hence \(f\) is bounded, and the lower and upper Riemann sums and integrals of \(f\) are defined.

Step 2: computing the sums for an equally spaced partition. Fix \(n \in \mathbb{Z}^+\) and let \(P_n\) be the partition \(x_0, x_1, \dots, x_n\) of \([a,b]\) given by

\begin{equation*} x_j = a + \frac{j(b-a)}{n}, \qquad j = 0, 1, \dots, n, \end{equation*}

so \(x_j - x_{j-1} = \frac{b-a}{n}\) for each \(j\).

Fix \(j \in \{1,\dots,n\}\). If \(x \in [x_{j-1},x_j]\), then monotonicity gives \(f(x_{j-1}) \le f(x) \le f(x_j)\) (this holds also when \(x\) equals one of the endpoints). Since \(x_{j-1}\) and \(x_j\) themselves belong to \([x_{j-1},x_j]\), these bounds are attained, so

\begin{equation*} \inf_{[x_{j-1},x_j]} f = f(x_{j-1}) \qquad \text{and} \qquad \sup_{[x_{j-1},x_j]} f = f(x_j). \end{equation*}

Therefore, by 1.3,

\begin{equation*} U(f,P_n,[a,b]) - L(f,P_n,[a,b]) = \frac{b-a}{n} \sum_{j=1}^{n} \big( f(x_j) - f(x_{j-1}) \big). \end{equation*}

The sum telescopes: \(\sum_{j=1}^{n}\big(f(x_j) - f(x_{j-1})\big) = f(x_n) - f(x_0) = f(b) - f(a)\). Hence

\begin{equation*} U(f,P_n,[a,b]) - L(f,P_n,[a,b]) = \frac{(b-a)\big(f(b) - f(a)\big)}{n}. \end{equation*}

Step 3: conclusion. Let \(\varepsilon > 0\). The number \((b-a)(f(b)-f(a))\) is a fixed nonnegative real number (nonnegative since \(a < b\) implies \(f(a) \le f(b)\)), so by the Archimedean property of \(\mathbb{R}\) there exists \(n \in \mathbb{Z}^+\) with

\begin{equation*} \frac{(b-a)\big(f(b)-f(a)\big)}{n} < \varepsilon . \end{equation*}

For this \(n\), Step 2 gives \(U(f,P_n,[a,b]) - L(f,P_n,[a,b]) < \varepsilon\).

Thus for every \(\varepsilon > 0\) there is a partition of \([a,b]\) whose upper and lower Riemann sums differ by less than \(\varepsilon\). By the Riemann criterion (Exercise 3 in Section 1A, proved in the solution to Exercise 9), \(f\) is Riemann integrable on \([a,b]\). \(\blacksquare\)

Problem (1A.14)

Suppose \(f_1, f_2, \dots\) is a sequence of Riemann integrable functions on \([a,b]\) such that \(f_1, f_2, \dots\) converges uniformly on \([a,b]\) to a function \(f \colon [a,b] \to \mathbb{R}\). Prove that \(f\) is Riemann integrable and

\begin{equation*} \int_a^b f = \lim_{n \to \infty} \int_a^b f_n . \end{equation*}

Solution

For \(n \in \mathbb{Z}^+\) set

\begin{equation*} \varepsilon_n = \sup_{x \in [a,b]} |f_n(x) - f(x)| \in [0,\infty]. \end{equation*}

Uniform convergence says exactly that for each \(\eta > 0\) there is \(N\) with \(|f_n(x) - f(x)| < \eta\) for all \(n \ge N\) and all \(x \in [a,b]\); hence \(\varepsilon_n \le \eta\) for all \(n \ge N\). In particular \(\varepsilon_n\) is finite for all large \(n\) and \(\varepsilon_n \to 0\).

Step 1: \(f\) is bounded. Choose \(N\) with \(\varepsilon_N \le 1\). Since \(f_N\) is Riemann integrable it is bounded, say \(|f_N| \le M\) on \([a,b]\). Then for every \(x \in [a,b]\),

\begin{equation*} |f(x)| \le |f(x) - f_N(x)| + |f_N(x)| \le 1 + M . \end{equation*}

So \(f\) is bounded and its lower and upper Riemann sums and integrals are defined. Note also that each \(f_n\) is bounded (being Riemann integrable), so from now on \(\varepsilon_n < \infty\) for every \(n \in \mathbb{Z}^+\).

Step 2: an oscillation comparison. Let \(A \subseteq [a,b]\) be nonempty and let \(n\) be such that \(\varepsilon_n < \infty\). For every \(x \in A\) we have \(f(x) \le f_n(x) + \varepsilon_n\) and \(f(x) \ge f_n(x) - \varepsilon_n\), hence

\begin{equation*} \sup_A f \le \sup_A f_n + \varepsilon_n \qquad \text{and} \qquad \inf_A f \ge \inf_A f_n - \varepsilon_n . \end{equation*}

Subtracting,

\begin{equation*} \sup_A f - \inf_A f \;\le\; \Big( \sup_A f_n - \inf_A f_n \Big) + 2\varepsilon_n . \end{equation*}

Applying this with \(A = [x_{j-1},x_j]\) for each subinterval of a partition \(P\) of \([a,b]\), multiplying by \(x_j - x_{j-1} > 0\) and summing (the lengths sum to \(b-a\)), we obtain by 1.3

\begin{equation*} U(f,P,[a,b]) - L(f,P,[a,b]) \;\le\; \Big( U(f_n,P,[a,b]) - L(f_n,P,[a,b]) \Big) + 2(b-a)\varepsilon_n . \qquad (\star) \end{equation*}

Step 3: \(f\) is Riemann integrable. Let \(\eta > 0\). Since \(\varepsilon_n \to 0\), choose \(n\) with

\begin{equation*} \varepsilon_n < \frac{\eta}{4(b-a)} . \end{equation*}

Since \(f_n\) is Riemann integrable, the Riemann criterion (Exercise 3 in Section 1A, proved in the solution to Exercise 9) provides a partition \(P\) of \([a,b]\) with

\begin{equation*} U(f_n,P,[a,b]) - L(f_n,P,[a,b]) < \frac{\eta}{2} . \end{equation*}

Then \((\star)\) gives

\begin{equation*} U(f,P,[a,b]) - L(f,P,[a,b]) < \frac{\eta}{2} + 2(b-a) \cdot \frac{\eta}{4(b-a)} = \eta . \end{equation*}

As \(\eta > 0\) was arbitrary, the same criterion shows that \(f\) is Riemann integrable on \([a,b]\).

Step 4: convergence of the integrals. Fix \(n\) with \(\varepsilon_n < \infty\). By Exercise 5 in Section 1A the function \(-f_n\) is Riemann integrable with \(\int_a^b (-f_n) = -\int_a^b f_n\), and by Step 3 the function \(f\) is Riemann integrable; so by Exercise 4 in Section 1A the function \(f - f_n\) is Riemann integrable with

\begin{equation*} \int_a^b (f - f_n) = \int_a^b f - \int_a^b f_n . \end{equation*}

By Exercise 12 in Section 1A the function \(|f - f_n|\) is Riemann integrable and

\begin{equation*} \left| \int_a^b f - \int_a^b f_n \right| = \left| \int_a^b (f - f_n) \right| \le \int_a^b |f - f_n| . \end{equation*}

Finally 1.13, applied to the Riemann integrable function \(|f-f_n|\), gives

\begin{equation*} \int_a^b |f - f_n| \le (b-a) \sup_{[a,b]} |f - f_n| = (b-a)\varepsilon_n . \end{equation*}

Combining the last two displays,

\begin{equation*} \left| \int_a^b f - \int_a^b f_n \right| \le (b-a) \varepsilon_n \qquad \text{whenever } \varepsilon_n < \infty . \end{equation*}

Now let \(\eta > 0\). Since \(\varepsilon_n \to 0\), there is \(N \in \mathbb{Z}^+\) such that \(\varepsilon_n < \dfrac{\eta}{b-a}\) for all \(n \ge N\) (in particular \(\varepsilon_n < \infty\) for such \(n\)), and then

\begin{equation*} \left| \int_a^b f - \int_a^b f_n \right| \le (b-a)\varepsilon_n < \eta \qquad \text{for all } n \ge N . \end{equation*}

Hence

\begin{equation*} \int_a^b f = \lim_{n \to \infty} \int_a^b f_n . \qquad \blacksquare \end{equation*}

Exercises 1B

Problem (1B.1)

Define \(f \colon [0,1] \to \mathbb{R}\) as follows:

\begin{equation*} f(a) = \begin{cases} 0 & \text{if } a \text{ is irrational,} \\ \dfrac{1}{n} & \text{if } a \text{ is rational and } n \text{ is the smallest positive} \\ & \text{integer such that } a = \frac{m}{n} \text{ for some integer } m. \end{cases} \end{equation*}

Show that \(f\) is Riemann integrable and compute \(\int_0^1 f\).

Solution

We show that \(f\) is Riemann integrable with \(\int_0^1 f = 0\).

Step 0: \(f\) is bounded, with \(0 \le f \le 1\).

If \(a\) is irrational then \(f(a) = 0\). If \(a\) is rational, then the smallest positive integer \(n\) with \(a = m/n\) for some integer \(m\) satisfies \(n \ge 1\), so \(f(a) = 1/n \in (0,1]\). Hence \(0 \le f(a) \le 1\) for all \(a \in [0,1]\), so \(f\) is a bounded function on \([0,1]\) and the lower and upper Riemann integrals of 1.7 are defined.

Step 1: \(L(f,[0,1]) = 0\).

Let \(P\) be the partition \(0 = x_0 < x_1 < \cdots < x_n = 1\) of \([0,1]\). Each subinterval \([x_{j-1}, x_j]\) has positive length and therefore contains an irrational number, at which \(f\) takes the value \(0\). Since \(f \ge 0\), this gives

\begin{equation*} \inf_{[x_{j-1},x_j]} f = 0 \qquad \text{for each } j = 1, \dots, n . \end{equation*}

Hence \(L(f, P, [0,1]) = \sum_{j=1}^{n} (x_j - x_{j-1}) \cdot 0 = 0\). As \(P\) was an arbitrary partition, taking the supremum over all \(P\) in 1.7 gives \(L(f,[0,1]) = 0\).

Step 2: the level sets of \(f\) are finite.

For \(N \in \mathbb{Z}^+\) put

\begin{equation*} A_N = \Bigl\{ a \in [0,1] : f(a) \ge \tfrac{1}{N} \Bigr\} . \end{equation*}

Suppose \(a \in A_N\). Then \(f(a) > 0\), so \(a\) is rational; writing \(n\) for the smallest positive integer with \(a = m/n\) for some integer \(m\), we have \(1/n = f(a) \ge 1/N\), hence \(n \le N\). Since \(a \in [0,1]\) and \(a = m/n\) with \(n > 0\), we get \(0 \le m \le n\). Therefore

\begin{equation*} A_N \subseteq \Bigl\{ \tfrac{m}{n} : n \in \{1, \dots, N\},\ m \in \{0, 1, \dots, n\} \Bigr\}, \end{equation*}

which is a finite set. Let \(K_N\) denote the number of elements of \(A_N\); note \(K_N \ge 1\) because \(0 = \frac{0}{1}\) lies in \(A_N\).

Moreover, if \(x \in [0,1] \setminus A_N\) then \(f(x) < 1/N\); since every value of \(f\) lies in \(\{0\} \cup \{1/n : n \in \mathbb{Z}^+\}\), this forces \(f(x) \le \frac{1}{N+1}\).

Step 3: \(U(f,[0,1]) = 0\).

Let \(\varepsilon > 0\). Choose \(N \in \mathbb{Z}^+\) with \(\frac{1}{N} < \frac{\varepsilon}{2}\), and set \(K = K_N\) as in Step 2. Now choose \(n \in \mathbb{Z}^+\) with

\begin{equation*} n > \frac{4K}{\varepsilon}, \end{equation*}

and let \(P_n\) be the equally spaced partition \(0, \frac{1}{n}, \frac{2}{n}, \dots, 1\) of \([0,1]\), whose subintervals \(I_j = \bigl[\frac{j-1}{n}, \frac{j}{n}\bigr]\) each have length \(\frac{1}{n}\).

Split the index set \(\{1, \dots, n\}\) into

\begin{equation*} J = \{ j : I_j \cap A_N \ne \emptyset \}, \qquad J^{\prime} = \{1,\dots,n\} \setminus J . \end{equation*}

Each point of \([0,1]\) lies in at most two of the intervals \(I_1, \dots, I_n\) (in two only when it is one of the interior partition points), so each of the \(K\) points of \(A_N\) contributes at most two indices to \(J\). Hence \(J\) has at most \(2K\) elements.

For \(j \in J\) we use only the crude bound \(\sup_{I_j} f \le 1\) from Step 0. For \(j \in J^{\prime}\) the interval \(I_j\) misses \(A_N\), so by the last sentence of Step 2 every \(x \in I_j\) satisfies \(f(x) \le \frac{1}{N+1} < \frac{1}{N}\), whence \(\sup_{I_j} f \le \frac{1}{N}\). Therefore

\begin{equation*} \begin{aligned} U(f, P_n, [0,1]) &= \sum_{j \in J} \frac{1}{n} \sup_{I_j} f \; + \; \sum_{j \in J^{\prime}} \frac{1}{n} \sup_{I_j} f \\ &\le \frac{2K}{n} \cdot 1 \; + \; \Bigl( \sum_{j \in J^{\prime}} \frac{1}{n} \Bigr) \frac{1}{N} \\ &\le \frac{2K}{n} + \frac{1}{N} \\ &< \frac{\varepsilon}{2} + \frac{\varepsilon}{2} = \varepsilon , \end{aligned} \end{equation*}

where we used \(\sum_{j \in J^{\prime}} \frac{1}{n} \le 1\) together with the choices \(n > 4K/\varepsilon\) and \(1/N < \varepsilon/2\).

By the definition 1.7 of the upper Riemann integral as an infimum over partitions, \(U(f,[0,1]) \le U(f,P_n,[0,1]) < \varepsilon\). Since \(\varepsilon > 0\) was arbitrary, \(U(f,[0,1]) \le 0\). On the other hand \(L(f,[0,1]) \le U(f,[0,1])\) by 1.8, and \(L(f,[0,1]) = 0\) by Step 1, so \(U(f,[0,1]) \ge 0\).

Hence \(U(f,[0,1]) = 0 = L(f,[0,1])\), so by 1.9 the function \(f\) is Riemann integrable and

\begin{equation*} \int_0^1 f = 0 . \end{equation*}

It is worth contrasting this with 1.14. The function here is discontinuous at every rational number of \([0,1]\), a dense infinite set, yet it is still Riemann integrable, because at each threshold \(1/N\) the set of points where \(f\) is at least that large is only finite.

Problem (1B.2)

Suppose \(f \colon [a,b] \to \mathbb{R}\) is a bounded function. Prove that \(f\) is Riemann integrable if and only if

\begin{equation*} L(-f, [a,b]) = -L(f, [a,b]) . \end{equation*}

Solution

The whole exercise rests on one identity, which says that passing to \(-f\) interchanges the roles of the lower and upper Riemann integrals:

\begin{equation*} U(f,[a,b]) = -L(-f,[a,b]) . \tag{\ast} \end{equation*}

First a remark about suprema and infima. If \(S\) is a nonempty set of real numbers that is bounded, then

\begin{equation*} \sup S = -\inf(-S), \qquad \text{where } -S = \{-s : s \in S\} . \end{equation*}

Indeed, \(\alpha\) is an upper bound for \(S\) exactly when \(-\alpha\) is a lower bound for \(-S\); hence the least upper bound of \(S\) is the negative of the greatest lower bound of \(-S\). Applying this with \(S\) the set of values of \(f\) on a nonempty subset \(E \subseteq [a,b]\), and noting that the set of values of \(-f\) on \(E\) is exactly \(-S\), we obtain

\begin{equation*} \sup_{E} f = -\inf_{E} (-f) . \tag{1} \end{equation*}

Note that \(-f\) is bounded because \(f\) is, so all quantities below are real numbers.

Now let \(P\) be any partition \(a = x_0 < x_1 < \cdots < x_n = b\) of \([a,b]\). Using the definitions of the upper and lower Riemann sums together with (1) applied on each \([x_{j-1},x_j]\),

\begin{equation*} \begin{aligned} U(f, P, [a,b]) &= \sum_{j=1}^{n} (x_j - x_{j-1}) \sup_{[x_{j-1},x_j]} f \\ &= \sum_{j=1}^{n} (x_j - x_{j-1}) \Bigl( -\inf_{[x_{j-1},x_j]} (-f) \Bigr) \\ &= -\sum_{j=1}^{n} (x_j - x_{j-1}) \inf_{[x_{j-1},x_j]} (-f) \\ &= -L(-f, P, [a,b]) . \end{aligned} \tag{2} \end{equation*}

Taking the infimum over all partitions \(P\) and using 1.7 twice, together with the remark above (in the form \(\inf T = -\sup(-T)\) applied to the set \(T\) of all numbers \(U(f,P,[a,b])\)),

\begin{equation*} U(f,[a,b]) = \inf_{P} U(f,P,[a,b]) = \inf_{P} \bigl( -L(-f,P,[a,b]) \bigr) = -\sup_{P} L(-f,P,[a,b]) = -L(-f,[a,b]), \end{equation*}

which is exactly \((\ast)\). Here the set \(T\) is nonempty and bounded below (by 1.8 applied to \(f\), or simply because \(U(f,P,[a,b]) \ge (b-a)\inf_{[a,b]} f\) for every \(P\)), so the manipulation is legitimate.

With \((\ast)\) in hand the equivalence is immediate. By 1.9, \(f\) is Riemann integrable if and only if

\begin{equation*} L(f,[a,b]) = U(f,[a,b]) . \end{equation*}

Substituting \((\ast)\) into the right side, this holds if and only if

\begin{equation*} L(f,[a,b]) = -L(-f,[a,b]), \end{equation*}

and multiplying both sides by \(-1\) (a reversible operation) this holds if and only if

\begin{equation*} L(-f,[a,b]) = -L(f,[a,b]) , \end{equation*}

as desired.

Two remarks. First, the exercise is genuinely a statement about \(f\) alone: it says that Riemann integrability can be tested using only the lower Riemann integral, applied to \(f\) and to \(-f\).

Second, \((\ast)\) also recovers Exercise 5 in Section 1A. Suppose \(f\) is Riemann integrable, so \(L(f,[a,b]) = U(f,[a,b]) = \int_a^b f\). By the equivalence just proved, \(L(-f,[a,b]) = -\int_a^b f\). Applying \((\ast)\) with \(-f\) in place of \(f\) (and using \(-(-f) = f\)) gives \(U(-f,[a,b]) = -L(f,[a,b]) = -\int_a^b f\) as well. Hence \(L(-f,[a,b]) = U(-f,[a,b])\), so \(-f\) is Riemann integrable and \(\int_a^b (-f) = -\int_a^b f\).

Problem (1B.3)

Suppose \(f, g \colon [a,b] \to \mathbb{R}\) are bounded functions. Prove that

\begin{equation*} L(f, [a,b]) + L(g, [a,b]) \le L(f+g, [a,b]) \end{equation*}

and

\begin{equation*} U(f+g, [a,b]) \le U(f, [a,b]) + U(g, [a,b]) . \end{equation*}

Solution

Since \(f\) and \(g\) are bounded, so is \(f+g\), and all six quantities in the statement are real numbers.

Step 1: pointwise inequalities for infima and suprema.

Let \(E \subseteq [a,b]\) be nonempty. For every \(x \in E\) we have \(\inf_E f \le f(x)\) and \(\inf_E g \le g(x)\), so

\begin{equation*} \inf_E f + \inf_E g \le f(x) + g(x) = (f+g)(x) . \end{equation*}

Thus \(\inf_E f + \inf_E g\) is a lower bound for the set of values of \(f+g\) on \(E\), and since \(\inf_E (f+g)\) is the greatest such lower bound,

\begin{equation*} \inf_E f + \inf_E g \le \inf_E (f+g) . \tag{1} \end{equation*}

Symmetrically, for every \(x \in E\) we have \((f+g)(x) = f(x) + g(x) \le \sup_E f + \sup_E g\), so \(\sup_E f + \sup_E g\) is an upper bound for the values of \(f+g\) on \(E\), and therefore

\begin{equation*} \sup_E (f+g) \le \sup_E f + \sup_E g . \tag{2} \end{equation*}

Step 2: the inequalities for a single partition.

Let \(P\) be the partition \(a = x_0 < x_1 < \cdots < x_n = b\) of \([a,b]\). Applying (1) with \(E = [x_{j-1},x_j]\) and multiplying by the positive number \(x_j - x_{j-1}\), then summing over \(j\),

\begin{equation*} L(f,P,[a,b]) + L(g,P,[a,b]) = \sum_{j=1}^{n} (x_j - x_{j-1}) \Bigl( \inf_{[x_{j-1},x_j]} f + \inf_{[x_{j-1},x_j]} g \Bigr) \le L(f+g, P, [a,b]) . \tag{3} \end{equation*}

Applying (2) in the same way,

\begin{equation*} U(f+g, P, [a,b]) \le U(f,P,[a,b]) + U(g,P,[a,b]) . \tag{4} \end{equation*}

Step 3: passing to the lower and upper Riemann integrals.

One cannot simply take a supremum in (3), because the suprema defining \(L(f,[a,b])\) and \(L(g,[a,b])\) may be approached along different partitions. The fix is to refine, using 1.5.

Let \(\varepsilon > 0\). By 1.7, \(L(f,[a,b])\) is the supremum of the lower Riemann sums of \(f\), so there is a partition \(P_1\) of \([a,b]\) with

\begin{equation*} L(f, P_1, [a,b]) > L(f,[a,b]) - \frac{\varepsilon}{2}, \end{equation*}

and likewise a partition \(P_2\) of \([a,b]\) with

\begin{equation*} L(g, P_2, [a,b]) > L(g,[a,b]) - \frac{\varepsilon}{2} . \end{equation*}

Let \(P\) be the partition of \([a,b]\) obtained by merging the lists that define \(P_1\) and \(P_2\). Then the list defining \(P_1\) is a sublist of the list defining \(P\), and the same holds for \(P_2\), so 1.5 gives

\begin{equation*} L(f, P, [a,b]) \ge L(f, P_1, [a,b]) \qquad \text{and} \qquad L(g, P, [a,b]) \ge L(g, P_2, [a,b]) . \end{equation*}

Combining these with (3) and with \(L(f+g,P,[a,b]) \le L(f+g,[a,b])\) (which holds by 1.7, as \(L(f+g,[a,b])\) is a supremum over all partitions),

\begin{equation*} \begin{aligned} L(f+g,[a,b]) &\ge L(f+g, P, [a,b]) \\ &\ge L(f, P, [a,b]) + L(g, P, [a,b]) \\ &\ge L(f, P_1, [a,b]) + L(g, P_2, [a,b]) \\ &> L(f,[a,b]) + L(g,[a,b]) - \varepsilon . \end{aligned} \end{equation*}

Since \(\varepsilon > 0\) was arbitrary, \(L(f,[a,b]) + L(g,[a,b]) \le L(f+g,[a,b])\), which is the first assertion.

The second assertion is proved by the mirror-image argument. Let \(\varepsilon > 0\). By 1.7 there are partitions \(Q_1, Q_2\) of \([a,b]\) with

\begin{equation*} U(f, Q_1, [a,b]) < U(f,[a,b]) + \frac{\varepsilon}{2} \qquad \text{and} \qquad U(g, Q_2, [a,b]) < U(g,[a,b]) + \frac{\varepsilon}{2} . \end{equation*}

Let \(Q\) be the partition obtained by merging the lists defining \(Q_1\) and \(Q_2\). By 1.5, adjoining points decreases upper Riemann sums, so \(U(f,Q,[a,b]) \le U(f,Q_1,[a,b])\) and \(U(g,Q,[a,b]) \le U(g,Q_2,[a,b])\). Hence, using (4),

\begin{equation*} \begin{aligned} U(f+g,[a,b]) &\le U(f+g, Q, [a,b]) \\ &\le U(f, Q, [a,b]) + U(g, Q, [a,b]) \\ &\le U(f, Q_1, [a,b]) + U(g, Q_2, [a,b]) \\ &< U(f,[a,b]) + U(g,[a,b]) + \varepsilon . \end{aligned} \end{equation*}

Since \(\varepsilon > 0\) was arbitrary, \(U(f+g,[a,b]) \le U(f,[a,b]) + U(g,[a,b])\), completing the proof.

Combining the two inequalities just proved with 1.8 gives a quick proof of Exercise 4 in Section 1A: if \(f\) and \(g\) are Riemann integrable, then

\begin{equation*} \int_a^b f + \int_a^b g = L(f,[a,b]) + L(g,[a,b]) \le L(f+g,[a,b]) \le U(f+g,[a,b]) \le U(f,[a,b]) + U(g,[a,b]) = \int_a^b f + \int_a^b g , \end{equation*}

so all the inequalities are equalities; in particular \(L(f+g,[a,b]) = U(f+g,[a,b])\), so \(f+g\) is Riemann integrable with \(\int_a^b (f+g) = \int_a^b f + \int_a^b g\).

Exercise 1B.4 shows that both inequalities in this exercise can be strict, so this last argument really does require the integrability of \(f\) and \(g\).

Problem (1B.4)

Give an example of bounded functions \(f, g \colon [0,1] \to \mathbb{R}\) such that

\begin{equation*} L(f, [0,1]) + L(g, [0,1]) < L(f+g, [0,1]) \end{equation*}

and

\begin{equation*} U(f+g, [0,1]) < U(f, [0,1]) + U(g, [0,1]) . \end{equation*}

Solution

Take \(f\) to be the function of 1.14 and \(g = 1 - f\). Explicitly, define \(f, g \colon [0,1] \to \mathbb{R}\) by

\begin{equation*} f(x) = \begin{cases} 1 & \text{if } x \text{ is rational,} \\ 0 & \text{if } x \text{ is irrational,} \end{cases} \qquad g(x) = \begin{cases} 0 & \text{if } x \text{ is rational,} \\ 1 & \text{if } x \text{ is irrational.} \end{cases} \end{equation*}

Both take only the values \(0\) and \(1\), so both are bounded.

The computation in 1.14 shows that \(L(f,[0,1]) = 0\) and \(U(f,[0,1]) = 1\). We recall the reason. If \([c,d] \subseteq [0,1]\) with \(c < d\), then \([c,d]\) contains an irrational number and also contains a rational number, so

\begin{equation*} \inf_{[c,d]} f = 0 \qquad \text{and} \qquad \sup_{[c,d]} f = 1 . \end{equation*}

Hence for every partition \(P\) of \([0,1]\), say \(0 = x_0 < x_1 < \cdots < x_n = 1\),

\begin{equation*} L(f,P,[0,1]) = \sum_{j=1}^n (x_j - x_{j-1}) \cdot 0 = 0, \qquad U(f,P,[0,1]) = \sum_{j=1}^n (x_j - x_{j-1}) \cdot 1 = 1 . \end{equation*}

Taking the supremum and the infimum over all partitions as in 1.7 gives \(L(f,[0,1]) = 0\) and \(U(f,[0,1]) = 1\).

The same argument applies verbatim to \(g\), since every nonempty subinterval of positive length also contains both a rational and an irrational number and \(g\) likewise takes the values \(0\) and \(1\) on each such subinterval. Hence

\begin{equation*} L(g,[0,1]) = 0 \qquad \text{and} \qquad U(g,[0,1]) = 1 . \end{equation*}

Now \(f + g\) is the constant function \(1\) on \([0,1]\): if \(x\) is rational then \(f(x) + g(x) = 1 + 0 = 1\), and if \(x\) is irrational then \(f(x) + g(x) = 0 + 1 = 1\). For a constant function every lower and every upper Riemann sum equals \(1\), because \(\inf_{[c,d]}(f+g) = \sup_{[c,d]}(f+g) = 1\) on each subinterval and the subinterval lengths of a partition of \([0,1]\) sum to \(1\). Hence

\begin{equation*} L(f+g,[0,1]) = U(f+g,[0,1]) = 1 . \end{equation*}

(Equivalently: \(f+g\) is continuous, so it is Riemann integrable by 1.11, and \(\int_0^1 (f+g) = 1\) by 1.13, since \(\inf_{[0,1]}(f+g) = \sup_{[0,1]}(f+g) = 1\).)

Putting these together:

\begin{equation*} L(f,[0,1]) + L(g,[0,1]) = 0 + 0 = 0 < 1 = L(f+g,[0,1]) , \end{equation*}

and

\begin{equation*} U(f+g,[0,1]) = 1 < 2 = 1 + 1 = U(f,[0,1]) + U(g,[0,1]) . \end{equation*}

Both inequalities of Exercise 1B.3 are therefore strict for this pair, as required.

The mechanism is worth naming. The lower Riemann integral fails to be additive because \(f\) and \(g\) take their small values at different points: on every subinterval, \(f\) dips to \(0\) at the irrationals while \(g\) dips to \(0\) at the rationals, and no single point makes both dip at once. Dually for the suprema. Cancellation of this kind is invisible to the one-function-at-a-time bookkeeping of Riemann sums, and it is precisely the phenomenon that forces the hypothesis of Riemann integrability into Exercise 4 in Section 1A.

Problem (1B.5)

Give an example of a sequence of continuous real-valued functions \(f_1, f_2, \dots\) on \([0,1]\) and a continuous real-valued function \(f\) on \([0,1]\) such that

\begin{equation*} f(x) = \lim_{k \to \infty} f_k(x) \end{equation*}

for each \(x \in [0,1]\) but

\begin{equation*} \int_0^1 f \ne \lim_{k \to \infty} \int_0^1 f_k . \end{equation*}

Solution

Take \(f\) to be the constant function \(0\), and let \(f_k\) be a tall thin triangular spike of area \(1\) sitting on the interval \([0, \frac1k]\). Explicitly, for \(k \in \mathbb{Z}^+\) define \(f_k \colon [0,1] \to \mathbb{R}\) by

\begin{equation*} f_k(x) = \begin{cases} 4k^2 x & \text{if } 0 \le x \le \frac{1}{2k}, \\[2pt] 4k^2 \bigl( \frac{1}{k} - x \bigr) & \text{if } \frac{1}{2k} \le x \le \frac{1}{k}, \\[2pt] 0 & \text{if } \frac{1}{k} \le x \le 1, \end{cases} \end{equation*}

and define \(f \colon [0,1] \to \mathbb{R}\) by \(f(x) = 0\).

Each \(f_k\) is well defined and continuous. The three formulas agree at the two points where the cases overlap: at \(x = \frac{1}{2k}\) the first gives \(4k^2 \cdot \frac{1}{2k} = 2k\) and the second gives \(4k^2 \bigl( \frac1k - \frac{1}{2k} \bigr) = 4k^2 \cdot \frac{1}{2k} = 2k\); at \(x = \frac1k\) the second gives \(0\) and the third gives \(0\). Each piece is a polynomial and hence continuous on its closed subinterval, so \(f_k\) is continuous on \([0,1]\) (this uses that the three subintervals are closed and cover \([0,1]\), with agreement at the overlaps). Also note \(\frac1k \le 1\), so the description makes sense for every \(k \in \mathbb{Z}^+\). The function \(f\) is constant, hence continuous. By 1.11, \(f\) and every \(f_k\) are Riemann integrable on \([0,1]\).

Pointwise convergence. Let \(x \in [0,1]\). If \(x = 0\), then \(f_k(0) = 4k^2 \cdot 0 = 0\) for every \(k\), so \(\lim_{k \to \infty} f_k(0) = 0 = f(0)\). If \(x > 0\), choose \(K \in \mathbb{Z}^+\) with \(K > \frac1x\), which is possible by the Archimedean property. For every \(k \ge K\) we have \(\frac1k \le \frac1K < x \le 1\), so \(x\) lies in the third case and \(f_k(x) = 0\). Hence \(\lim_{k \to \infty} f_k(x) = 0 = f(x)\). Thus \(f(x) = \lim_{k \to \infty} f_k(x)\) for each \(x \in [0,1]\), as required.

The integrals. Clearly \(\int_0^1 f = 0\), since \(\inf_{[0,1]} f = \sup_{[0,1]} f = 0\) and hence \(\int_0^1 f = 0\) by 1.13.

For \(f_k\), first suppose \(k \ge 2\), so that \(\frac1k \in (0,1)\) and \(\frac{1}{2k} \in (0, \frac1k)\). Apply Exercise 10 in Section 1A once to split \([0,1]\) at \(\frac1k\), and again to split \([0,\frac1k]\) at \(\frac{1}{2k}\) (each application is legitimate because \(f_k\) is Riemann integrable on \([0,1]\), hence on every closed subinterval by Exercise 9 in Section 1A):

\begin{equation*} \int_0^1 f_k = \int_0^{1/(2k)} f_k + \int_{1/(2k)}^{1/k} f_k + \int_{1/k}^{1} f_k . \end{equation*}

When \(k = 1\) we have \(\frac1k = 1\), so only the split at \(\frac{1}{2k} = \frac12\) is used and the identity above holds with its last term read as \(0\). On \([0, \frac{1}{2k}]\) the function \(f_k\) agrees with \(x \mapsto 4k^2 x\), whose antiderivative is \(2k^2 x^2\), so by the fundamental theorem of calculus

\begin{equation*} \int_0^{1/(2k)} f_k = 2k^2 \Bigl( \frac{1}{2k} \Bigr)^{\!2} = 2k^2 \cdot \frac{1}{4k^2} = \frac12 . \end{equation*}

On \([\frac{1}{2k}, \frac1k]\) the function \(f_k\) agrees with \(x \mapsto 4k^2 (\frac1k - x)\), whose antiderivative is \(-2k^2 (\frac1k - x)^2\), so

\begin{equation*} \int_{1/(2k)}^{1/k} f_k = \Bigl[ -2k^2 \Bigl( \tfrac1k - x \Bigr)^{\!2} \Bigr]_{x = 1/(2k)}^{x = 1/k} = 0 - \Bigl( -2k^2 \cdot \frac{1}{4k^2} \Bigr) = \frac12 . \end{equation*}

On \([\frac1k, 1]\) the function \(f_k\) is identically \(0\), so that integral is \(0\) by 1.13. Adding,

\begin{equation*} \int_0^1 f_k = \frac12 + \frac12 + 0 = 1 \qquad \text{for every } k \in \mathbb{Z}^+ . \end{equation*}

Geometrically this is just the area of a triangle with base \(\frac1k\) and height \(2k\), namely \(\frac12 \cdot \frac1k \cdot 2k = 1\).

Conclusion. The limit \(\lim_{k \to \infty} \int_0^1 f_k = \lim_{k \to \infty} 1 = 1\) exists, but

\begin{equation*} \int_0^1 f = 0 \ne 1 = \lim_{k \to \infty} \int_0^1 f_k , \end{equation*}

so this sequence has the required property.

Note that this example does not contradict 1.18. There the sequence is assumed uniformly bounded, that is, \(|f_k(x)| \le M\) for all \(k\) and all \(x\) with a single constant \(M\). Here \(\sup_{[0,1]} f_k = f_k \bigl( \frac{1}{2k} \bigr) = 2k \to \infty\), so no such \(M\) exists. The example shows that the uniform boundedness hypothesis in 1.18 cannot be dropped, and it complements 1.17: there the failure was that the pointwise limit need not be Riemann integrable, whereas here the limit function is as well behaved as could be wished and the integrals still fail to converge to the right value. The mass of \(f_k\) escapes by concentrating on shorter and shorter intervals rather than by growing in the limit function.

Measures

Exercises 2A

Problem (2A.1)

Prove that if \(A\) and \(B\) are subsets of \(\mathbb{R}\) and \(|B| = 0\), then \(|A \cup B| = |A|\).

Solution

We prove the two inequalities separately.

First, \(A \subseteq A \cup B\). Because outer measure preserves order (2.5), this gives

\begin{equation*} |A| \le |A \cup B|. \end{equation*}

Second, apply countable subadditivity of outer measure (2.8) to the sequence \(A, B, \emptyset, \emptyset, \dots\), whose union is \(A \cup B\). Since \(|\emptyset| = 0\) (the empty set is finite, so 2.3 applies, or take every \(I_k = \emptyset\) in 2.2), we obtain

\begin{equation*} |A \cup B| \le |A| + |B| = |A| + 0 = |A|. \end{equation*}

Combining the two displayed inequalities yields \(|A \cup B| = |A|\). Note that the argument is valid also when \(|A| = \infty\), since both inequalities are statements in \([0, \infty]\) and no subtraction was performed.

Problem (2A.2)

Suppose \(A \subseteq \mathbb{R}\) and \(t \in \mathbb{R}\). Let \(tA = \{ta : a \in A\}\). Prove that \(|tA| = |t|\,|A|\).

[Assume that \(0 \cdot \infty\) is defined to be \(0\).]

Solution

Case 1: \(t = 0\). If \(A = \emptyset\) then \(tA = \emptyset\), and if \(A \ne \emptyset\) then \(tA = \{0\}\). In either case \(tA\) is a finite set, so \(|tA| = 0\) by 2.3. On the other side, \(|t|\,|A| = 0 \cdot |A| = 0\), using the stated convention \(0 \cdot \infty = 0\) in case \(|A| = \infty\). Hence \(|tA| = |t|\,|A|\).

Case 2: \(t \ne 0\). We first record how dilation acts on open intervals: for every open interval \(I\),

\begin{equation*} \ell(tI) = |t|\,\ell(I), \end{equation*}

where \(tI = \{tx : x \in I\}\) is again an open interval. Indeed, if \(I = \emptyset\) then \(tI = \emptyset\) and both sides are \(0\). If \(I = (\alpha, \beta)\) with \(\alpha, \beta \in \mathbb{R}\) and \(\alpha < \beta\), then \(tI = (t\alpha, t\beta)\) when \(t > 0\) and \(tI = (t\beta, t\alpha)\) when \(t < 0\); in both cases \(\ell(tI) = |t|(\beta - \alpha) = |t|\,\ell(I)\). If \(I\) is one of \((-\infty, \alpha)\), \((\alpha, \infty)\), \((-\infty, \infty)\), then \(tI\) is again an unbounded open interval (multiplication by \(t \ne 0\) is a homeomorphism of \(\mathbb{R}\) carrying rays to rays), so \(\ell(tI) = \infty = |t| \cdot \infty = |t|\,\ell(I)\), the middle equality because \(|t| > 0\).

Now we prove \(|tA| \le |t|\,|A|\). If \(|A| = \infty\) this is immediate, so assume \(|A| < \infty\) and let \(I_1, I_2, \dots\) be any sequence of open intervals with \(A \subseteq \bigcup_{k=1}^{\infty} I_k\). Applying the map \(x \mapsto tx\), which is a bijection of \(\mathbb{R}\), we get \(tA \subseteq \bigcup_{k=1}^{\infty} t I_k\), and each \(tI_k\) is an open interval. Hence by the definition 2.2 of outer measure,

\begin{equation*} |tA| \le \sum_{k=1}^{\infty} \ell(t I_k) = |t| \sum_{k=1}^{\infty} \ell(I_k). \end{equation*}

Taking the infimum of the right-hand side over all such sequences \(I_1, I_2, \dots\) gives \(|tA| \le |t|\,|A|\).

Finally we get the reverse inequality by applying the inequality just proved with \(t\) replaced by \(1/t\) and \(A\) replaced by \(tA\). Since \((1/t)(tA) = A\), this yields

\begin{equation*} |A| \le \tfrac{1}{|t|}\,|tA|, \end{equation*}

that is, \(|t|\,|A| \le |tA|\). Combining the two inequalities gives \(|tA| = |t|\,|A|\).

Problem (2A.3)

Prove that if \(A, B \subseteq \mathbb{R}\) and \(|A| < \infty\), then \(|B \setminus A| \ge |B| - |A|\).

Solution

Every element of \(B\) either lies in \(A\) or does not, so

\begin{equation*} B \subseteq (B \setminus A) \cup A. \end{equation*}

Because outer measure preserves order (2.5) and is countably subadditive (2.8, applied to the sequence \(B \setminus A,\ A,\ \emptyset,\ \emptyset, \dots\), which gives finite subadditivity), we obtain

\begin{equation*} |B| \le \bigl| (B \setminus A) \cup A \bigr| \le |B \setminus A| + |A|. \end{equation*}

If \(|B \setminus A| = \infty\), the desired inequality \(|B \setminus A| \ge |B| - |A|\) holds trivially. Otherwise \(|B \setminus A| < \infty\), and since also \(|A| < \infty\) by hypothesis, both terms on the right side of the display above are real numbers; hence \(|B| < \infty\) as well and we may subtract \(|A|\) from both sides of the inequality \(|B| \le |B \setminus A| + |A|\), obtaining

\begin{equation*} |B| - |A| \le |B \setminus A|, \end{equation*}

as desired. Note that the hypothesis \(|A| < \infty\) is what makes the subtraction legitimate; without it the expression \(|B| - |A|\) could be the undefined \(\infty - \infty\).

Problem (2A.4)

Suppose \(F\) is a subset of \(\mathbb{R}\) with the property that every open cover of \(F\) has a finite subcover. Prove that \(F\) is closed and bounded.

Solution

Together with the Heine–Borel Theorem (2.12), this exercise shows that a subset of \(\mathbb{R}\) has the finite-subcover property if and only if it is closed and bounded.

\(F\) is bounded. Consider the collection

\begin{equation*} \mathcal{C} = \{ (-k, k) : k \in \mathbb{Z}^{+} \} \end{equation*}

of open subsets of \(\mathbb{R}\). Every real number \(x\) satisfies \(|x| < k\) for some \(k \in \mathbb{Z}^{+}\), so \(\bigcup_{k=1}^{\infty} (-k,k) = \mathbb{R} \supseteq F\); thus \(\mathcal{C}\) is an open cover of \(F\) in the sense of 2.10. By hypothesis there exist \(k_1, \dots, k_n \in \mathbb{Z}^{+}\) with

\begin{equation*} F \subseteq (-k_1, k_1) \cup \cdots \cup (-k_n, k_n). \end{equation*}

(If the finite subcover is the empty list, then \(F = \emptyset\) and \(F\) is certainly bounded.) The intervals \((-k,k)\) increase with \(k\), so setting \(K = \max\{k_1, \dots, k_n\}\) the union above equals \((-K, K)\). Hence \(|x| < K\) for every \(x \in F\), so \(F\) is bounded.

\(F\) is closed. We show that the complement \(\mathbb{R} \setminus F\) is open. Let \(b \in \mathbb{R} \setminus F\). For each \(k \in \mathbb{Z}^{+}\) define

\begin{equation*} G_k = \Bigl(-\infty,\ b - \tfrac{1}{k}\Bigr) \cup \Bigl(b + \tfrac{1}{k},\ \infty\Bigr) = \Bigl\{ x \in \mathbb{R} : |x - b| > \tfrac{1}{k} \Bigr\}, \end{equation*}

which is an open subset of \(\mathbb{R}\), being a union of two open intervals. If \(|x - b| > 1/k\) then \(|x - b| > 1/(k+1)\), so \(G_1 \subseteq G_2 \subseteq G_3 \subseteq \cdots\). Moreover

\begin{equation*} \bigcup_{k=1}^{\infty} G_k = \{ x \in \mathbb{R} : x \ne b \} = \mathbb{R} \setminus \{b\}, \end{equation*}

because any \(x \ne b\) satisfies \(|x - b| > 1/k\) for all sufficiently large \(k\). Since \(b \notin F\), we have \(F \subseteq \mathbb{R} \setminus \{b\}\), so \(\{G_k : k \in \mathbb{Z}^{+}\}\) is an open cover of \(F\).

By hypothesis this cover has a finite subcover: there exist \(k_1, \dots, k_n \in \mathbb{Z}^{+}\) with \(F \subseteq G_{k_1} \cup \cdots \cup G_{k_n}\). (Again, if the list is empty then \(F = \emptyset\), which is closed, and we are done.) Because the \(G_k\) are nested increasing, with \(K = \max\{k_1, \dots, k_n\}\) we get \(F \subseteq G_K\). Thus no point of \(F\) lies within distance \(1/K\) of \(b\), that is,

\begin{equation*} \Bigl( b - \tfrac{1}{K},\ b + \tfrac{1}{K} \Bigr) \subseteq \mathbb{R} \setminus F . \end{equation*}

So every point of \(\mathbb{R} \setminus F\) is contained in an open interval contained in \(\mathbb{R} \setminus F\), which means \(\mathbb{R} \setminus F\) is open. Hence \(F\) is closed.

Problem (2A.5)

Suppose \(\mathcal{A}\) is a set of closed subsets of \(\mathbb{R}\) such that \(\bigcap_{F \in \mathcal{A}} F = \emptyset\). Prove that if \(\mathcal{A}\) contains at least one bounded set, then there exist \(n \in \mathbb{Z}^{+}\) and \(F_1, \dots, F_n \in \mathcal{A}\) such that \(F_1 \cap \cdots \cap F_n = \emptyset\).

Solution

Let \(F_0 \in \mathcal{A}\) be a bounded set. Since every element of \(\mathcal{A}\) is closed, \(F_0\) is a closed bounded subset of \(\mathbb{R}\).

If \(F_0 = \emptyset\), take \(n = 1\) and \(F_1 = F_0\); then \(F_1 = \emptyset\) and we are done. So assume \(F_0 \ne \emptyset\).

Consider the collection

\begin{equation*} \mathcal{C} = \{\, \mathbb{R} \setminus F : F \in \mathcal{A} \,\}. \end{equation*}

Each element of \(\mathcal{C}\) is an open subset of \(\mathbb{R}\), because each \(F \in \mathcal{A}\) is closed. We claim \(\mathcal{C}\) is an open cover of \(F_0\). Indeed, let \(x \in F_0\). Since \(\bigcap_{F \in \mathcal{A}} F = \emptyset\), the point \(x\) does not belong to every element of \(\mathcal{A}\); hence there is some \(F \in \mathcal{A}\) with \(x \notin F\), i.e. \(x \in \mathbb{R} \setminus F \in \mathcal{C}\). Thus \(F_0\) is contained in the union of the sets in \(\mathcal{C}\), proving the claim.

By the Heine–Borel Theorem (2.12), the open cover \(\mathcal{C}\) of the closed bounded set \(F_0\) has a finite subcover: there exist \(m \in \mathbb{Z}^{+}\) and \(F_1, \dots, F_m \in \mathcal{A}\) such that

\begin{equation*} F_0 \subseteq (\mathbb{R} \setminus F_1) \cup \cdots \cup (\mathbb{R} \setminus F_m). \end{equation*}

(The subcover cannot be an empty list, since \(F_0 \ne \emptyset\).) By De Morgan’s laws the right-hand side equals \(\mathbb{R} \setminus (F_1 \cap \cdots \cap F_m)\), so

\begin{equation*} F_0 \subseteq \mathbb{R} \setminus (F_1 \cap \cdots \cap F_m), \end{equation*}

which says precisely that

\begin{equation*} F_0 \cap F_1 \cap \cdots \cap F_m = \emptyset. \end{equation*}

Since \(F_0, F_1, \dots, F_m\) are all elements of \(\mathcal{A}\), relabelling this list as \(F_1, \dots, F_n\) with \(n = m + 1\) gives \(n \in \mathbb{Z}^{+}\) and \(F_1, \dots, F_n \in \mathcal{A}\) with \(F_1 \cap \cdots \cap F_n = \emptyset\), as required.

Remark: the hypothesis that some element of \(\mathcal{A}\) be bounded cannot be dropped. For example, \(\mathcal{A} = \{ [k, \infty) : k \in \mathbb{Z}^{+} \}\) consists of closed sets with \(\bigcap_{F \in \mathcal{A}} F = \emptyset\), yet every finite subfamily has nonempty intersection, since \([k_1, \infty) \cap \cdots \cap [k_n, \infty) = [\max\{k_1,\dots,k_n\}, \infty) \ne \emptyset\).

Problem (2A.6)

Prove that if \(a, b \in \mathbb{R}\) and \(a < b\), then

\begin{equation*} |(a,b)| = |[a,b)| = |(a,b]| = b - a. \end{equation*}

Solution

By 2.14 we know that \(|[a,b]| = b - a\). Each of the three sets in question differs from \([a,b]\) by at most the two-point set \(\{a,b\}\), and finite sets have outer measure \(0\) by 2.3. We can therefore obtain all three equalities from Exercise 1 in this section, which says that adjoining a set of outer measure \(0\) does not change outer measure.

The open interval. Since \([a,b] = (a,b) \cup \{a,b\}\) and \(|\{a,b\}| = 0\), Exercise 1 (with \(A = (a,b)\) and \(B = \{a,b\}\)) gives

\begin{equation*} |(a,b)| = \bigl| (a,b) \cup \{a,b\} \bigr| = |[a,b]| = b - a. \end{equation*}

The half-open interval \([a,b)\). Since \([a,b] = [a,b) \cup \{b\}\) and \(|\{b\}| = 0\), Exercise 1 gives

\begin{equation*} |[a,b)| = \bigl| [a,b) \cup \{b\} \bigr| = |[a,b]| = b - a. \end{equation*}

The half-open interval \((a,b]\). Since \([a,b] = (a,b] \cup \{a\}\) and \(|\{a\}| = 0\), Exercise 1 gives

\begin{equation*} |(a,b]| = \bigl| (a,b] \cup \{a\} \bigr| = |[a,b]| = b - a. \end{equation*}

Hence \(|(a,b)| = |[a,b)| = |(a,b]| = b - a\).

For the reader who prefers an argument not relying on Exercise 1, here is a direct proof for \((a,b)\); the two half-open cases then follow since \((a,b) \subseteq [a,b) \subseteq [a,b]\) and \((a,b) \subseteq (a,b] \subseteq [a,b]\), so that 2.5 squeezes their outer measures between \(|(a,b)|\) and \(|[a,b]|\). For the upper bound, \((a,b) \subseteq [a,b]\), so \(|(a,b)| \le |[a,b]| = b - a\) by 2.5 and 2.14. For the lower bound, let \(\varepsilon\) satisfy \(0 < \varepsilon < (b-a)/2\). Then \([a + \varepsilon, b - \varepsilon] \subseteq (a,b)\), so by 2.5 and 2.14,

\begin{equation*} |(a,b)| \ge \bigl| [a + \varepsilon,\ b - \varepsilon] \bigr| = (b - \varepsilon) - (a + \varepsilon) = b - a - 2\varepsilon . \end{equation*}

Because this holds for every such \(\varepsilon\), letting \(\varepsilon\) tend to \(0\) gives \(|(a,b)| \ge b - a\).

Problem (2A.7)

Suppose \(a, b, c, d\) are real numbers with \(a < b\) and \(c < d\). Prove that

\begin{equation*} |(a,b) \cup (c,d)| = (b-a) + (d-c) \quad \text{if and only if} \quad (a,b) \cap (c,d) = \emptyset . \end{equation*}

Solution

Throughout write \(E = (a,b) \cup (c,d)\). Note first that by Exercise 6 of this section, \(|(a,b)| = b - a\) and \(|(c,d)| = d - c\), and hence by finite subadditivity (a special case of 2.8, applied to the sequence \((a,b),\ (c,d),\ \emptyset,\ \emptyset, \dots\)) we always have

\begin{equation*} |E| \le (b-a) + (d-c) . \end{equation*}

So the content of the exercise is that equality holds exactly when the two intervals are disjoint.

We begin with an elementary observation about when two nonempty open intervals meet.

Claim. \((a,b) \cap (c,d) \ne \emptyset\) if and only if \(c < b\) and \(a < d\); and in that case

\begin{equation*} (a,b) \cap (c,d) = (u,v), \qquad u = \max\{a,c\}, \quad v = \min\{b,d\}, \end{equation*}

with \(u < v\).

Proof of claim. A real number \(x\) lies in \((a,b) \cap (c,d)\) exactly when \(a < x < b\) and \(c < x < d\), i.e. exactly when \(\max\{a,c\} < x < \min\{b,d\}\); so the intersection is the interval \((u,v)\), which is nonempty precisely when \(u < v\). Now \(u < v\) means \(\max\{a,c\} < \min\{b,d\}\), i.e. all four of \(a < b\), \(a < d\), \(c < b\), \(c < d\) hold. Two of these, \(a<b\) and \(c<d\), hold by hypothesis, so \(u < v\) is equivalent to the conjunction of \(c < b\) and \(a < d\). This proves the claim.

Suppose \((a,b) \cap (c,d) = \emptyset\). By the claim, \(b \le c\) or \(d \le a\). The two cases are symmetric under interchanging the pair \((a,b)\) with the pair \((c,d)\), which changes neither \(E\) nor the quantity \((b-a)+(d-c)\); so we may assume \(b \le c\).

We must prove \(|E| \ge (b-a) + (d-c)\); combined with the subadditivity bound displayed above this gives equality. Since \(a < b \le c < d\), we claim that

\begin{equation*} (a,d) \subseteq (a,b) \cup [b,c] \cup (c,d) = E \cup [b,c]. \end{equation*}

Indeed, let \(x \in (a,d)\). If \(x < b\) then \(a < x < b\), so \(x \in (a,b)\). If \(b \le x \le c\) then \(x \in [b,c]\). If \(x > c\) then \(c < x < d\), so \(x \in (c,d)\). These three cases are exhaustive, which proves the inclusion. (Note that the closed interval \([a,d]\) could not be used in place of \((a,d)\) here: neither \(a\) nor \(d\) belongs to \(E \cup [b,c]\).)

Now \(|[b,c]| = c - b\): this is 2.14 if \(b < c\), and if \(b = c\) then \([b,c] = \{b\}\) is a one-point set, so \(|[b,c]| = 0 = c - b\) by 2.3. Also \(|(a,d)| = d - a\) by Exercise 6 of this section, which applies because \(a < d\). Applying 2.5 and then finite subadditivity (2.8) to the displayed inclusion,

\begin{equation*} d - a = |(a,d)| \le |E| + |[b,c]| = |E| + (c - b). \end{equation*}

All quantities here are finite, so we may rearrange:

\begin{equation*} |E| \ge (d - a) - (c - b) = (b - a) + (d - c). \end{equation*}

Together with \(|E| \le (b-a)+(d-c)\) this gives \(|E| = (b-a) + (d-c)\).

Suppose \((a,b) \cap (c,d) \ne \emptyset\). We show \(|E| < (b-a) + (d-c)\), which in particular gives \(|E| \ne (b-a)+(d-c)\) and hence, by contraposition, completes the proof of the equivalence.

Put \(p = \min\{a,c\}\), \(q = \max\{b,d\}\), and let \(u = \max\{a,c\}\), \(v = \min\{b,d\}\) as in the claim, so \(u < v\). Since \(\{p,u\}\) and \(\{a,c\}\) agree as multisets, and likewise \(\{v,q\}\) and \(\{b,d\}\), at least one of the two intervals has left endpoint \(p\) and at least one has right endpoint \(q\). We first check that \(E = (p,q)\). The inclusion \(E \subseteq (p,q)\) is clear, since \((a,b) \subseteq (p,q)\) and \((c,d) \subseteq (p,q)\). Conversely, let \(x \in (p,q)\), so \(p < x < q\). If \(x \le u\) then, since \(p < x\) and \(\{p, u\} = \{a, c\}\), the point \(x\) lies in whichever of \((a,b)\), \((c,d)\) has left endpoint \(p\); indeed that interval is \((p, w)\) for some \(w \in \{b,d\}\) with \(w \ge v > u \ge x\), so \(p < x < w\) and \(x\) belongs to it. If \(x > u\) then, since \(x < q\) and \(\{v,q\} = \{b,d\}\), the point \(x\) lies in whichever of \((a,b)\), \((c,d)\) has right endpoint \(q\); that interval is \((z, q)\) for some \(z \in \{a,c\}\) with \(z \le u < x\), so \(z < x < q\) and \(x\) belongs to it. Hence \((p,q) \subseteq E\), and therefore \(E = (p,q)\).

Since \(p \le u < v \le q\) we have \(p < q\), so Exercise 6 gives

\begin{equation*} |E| = |(p,q)| = q - p . \end{equation*}

Finally, because \(\{p, u\} = \{a, c\}\) and \(\{v, q\} = \{b, d\}\) as multisets, we have \(p + u = a + c\) and \(v + q = b + d\). Therefore

\begin{equation*} (b - a) + (d - c) = (b + d) - (a + c) = (v + q) - (p + u) = (q - p) + (v - u) = |E| + (v - u). \end{equation*}

Since \(v - u > 0\) and \(|E| = q - p\) is finite, this gives

\begin{equation*} |E| < (b-a) + (d-c), \end{equation*}

as claimed. This completes both directions of the proof.

Problem (2A.8)

Prove that if \(A \subseteq \mathbb{R}\) and \(t > 0\), then

\begin{equation*} |A| = |A \cap (-t,t)| + \bigl|A \cap \bigl(\mathbb{R} \setminus (-t,t)\bigr)\bigr|. \end{equation*}

Solution

Write \(B = A \cap (-t,t)\) and \(C = A \cap \bigl(\mathbb{R}\setminus(-t,t)\bigr) = A \setminus (-t,t)\), so that \(A = B \cup C\).

The inequality \(|A| \le |B| + |C|\). This is immediate from the countable subadditivity of outer measure (2.8), applied to the sequence \(B, C, \emptyset, \emptyset, \dots\) whose union is \(A\).

The inequality \(|A| \ge |B| + |C|\). If \(|A| = \infty\) there is nothing to prove, so assume \(|A| < \infty\). Let \(\varepsilon > 0\). By the definition of outer measure (2.2) there is a sequence \(I_1, I_2, \dots\) of open intervals with

\begin{equation*} A \subseteq \bigcup_{k=1}^{\infty} I_k \qquad\text{and}\qquad \sum_{k=1}^{\infty} \ell(I_k) \le |A| + \varepsilon . \end{equation*}

Cut each \(I_k\) at the two points \(-t\) and \(t\): put

\begin{equation*} J_k = I_k \cap (-t,t), \qquad K_k = I_k \cap (-\infty,-t), \qquad L_k = I_k \cap (t,\infty). \end{equation*}

Each of \(J_k, K_k, L_k\) is an open interval (possibly empty), being the intersection of two open intervals. I claim that

\begin{equation*} \ell(J_k) + \ell(K_k) + \ell(L_k) = \ell(I_k) \qquad \text{for each } k . \end{equation*}

To see this, first note the one-point version: if \(I\) is an open interval and \(c \in \mathbb{R}\), then

\begin{equation*} \ell\bigl(I \cap (-\infty,c)\bigr) + \ell\bigl(I \cap (c,\infty)\bigr) = \ell(I). \end{equation*}

Indeed, if \(I = \emptyset\) both sides are \(0\). If \(I = (a,b)\) with \(a,b \in \mathbb{R}\), then for \(c \le a\) the two pieces are \(\emptyset\) and \((a,b)\), for \(c \ge b\) they are \((a,b)\) and \(\emptyset\), and for \(a < c < b\) they are \((a,c)\) and \((c,b)\), whose lengths sum to \((c-a)+(b-c) = b-a\); in every case the sum is \(b - a = \ell(I)\). If \(I\) is unbounded, then at least one of the two pieces is an unbounded open interval, so both sides equal \(\infty\).

Now write \(M_k = I_k \cap (-t,\infty)\), which is again an open interval. Applying the one-point version to \(I_k\) with \(c = -t\) gives \(\ell(K_k) + \ell(M_k) = \ell(I_k)\). Applying it to \(M_k\) with \(c = t\) gives \(\ell\bigl(M_k \cap (-\infty,t)\bigr) + \ell\bigl(M_k \cap (t,\infty)\bigr) = \ell(M_k)\); here \(M_k \cap (-\infty,t) = I_k \cap (-t,t) = J_k\), and, because \(t > -t\), \(M_k \cap (t,\infty) = I_k \cap (t,\infty) = L_k\). Hence \(\ell(J_k) + \ell(L_k) = \ell(M_k)\), and adding the two identities gives the claim.

Now \(B = A \cap (-t,t) \subseteq \bigcup_{k} J_k\), so by 2.2,

\begin{equation*} |B| \le \sum_{k=1}^{\infty} \ell(J_k). \end{equation*}

Also

\begin{equation*} C = A \setminus (-t,t) \subseteq \Bigl(\bigcup_{k=1}^{\infty} K_k\Bigr) \cup \Bigl(\bigcup_{k=1}^{\infty} L_k\Bigr) \cup \{-t,t\}, \end{equation*}

because a point of \(A\) outside \((-t,t)\) lies in some \(I_k\), and if it is neither \(-t\) nor \(t\) it then lies in \(K_k\) or in \(L_k\). The sequence \(K_1, L_1, K_2, L_2, \dots\) is a sequence of open intervals, so by 2.2 the union of its terms has outer measure at most \(\sum_{k}\bigl(\ell(K_k) + \ell(L_k)\bigr)\); since \(|\{-t,t\}| = 0\) by 2.3, subadditivity (2.8) gives

\begin{equation*} |C| \le \sum_{k=1}^{\infty} \bigl(\ell(K_k) + \ell(L_k)\bigr). \end{equation*}

Adding the two displayed bounds and using the claim,

\begin{equation*} |B| + |C| \le \sum_{k=1}^{\infty} \bigl(\ell(J_k) + \ell(K_k) + \ell(L_k)\bigr) = \sum_{k=1}^{\infty} \ell(I_k) \le |A| + \varepsilon . \end{equation*}

Because \(\varepsilon > 0\) is arbitrary, \(|B| + |C| \le |A|\).

Combining the two inequalities gives \(|A| = |A \cap (-t,t)| + \bigl|A \cap (\mathbb{R}\setminus(-t,t))\bigr|\).

Problem (2A.9)

Prove that \(|A| = \lim\limits_{t \to \infty} |A \cap (-t,t)|\) for all \(A \subseteq \mathbb{R}\).

Solution

First note that the limit exists in \([0,\infty]\). If \(0 < s < t\), then \(A \cap (-s,s) \subseteq A \cap (-t,t)\), so \(|A \cap (-s,s)| \le |A \cap (-t,t)|\) because outer measure preserves order (2.5). Thus \(t \mapsto |A \cap (-t,t)|\) is a nondecreasing function of \(t\) on \((0,\infty)\), and hence

\begin{equation*} \lim_{t \to \infty} |A \cap (-t,t)| = \sup_{t > 0} |A \cap (-t,t)| = \lim_{n \to \infty} |A \cap (-n,n)|, \end{equation*}

the last equality because the supremum over \(t > 0\) of a nondecreasing function is already attained as a limit along the positive integers. Call this common value \(L \in [0,\infty]\).

The inequality \(L \le |A|\). For every \(t > 0\) we have \(A \cap (-t,t) \subseteq A\), so \(|A \cap (-t,t)| \le |A|\) by 2.5. Taking the supremum over \(t > 0\) gives \(L \le |A|\).

The inequality \(|A| \le L\). Decompose \(A\) into the annuli cut out by the integers: set

\begin{equation*} C_1 = A \cap (-1,1) \qquad\text{and}\qquad C_n = A \cap \{x \in \mathbb{R} : n - 1 \le |x| < n\} \ \text{ for } n \ge 2 . \end{equation*}

Fix \(n \ge 2\) and apply Exercise 8 of Section 2A (proved above) to the set \(A \cap (-n,n)\) with \(t = n-1\), which is legitimate because \(n - 1 > 0\). Since \((-(n-1), n-1) \subseteq (-n,n)\) we have

\begin{equation*} \bigl(A \cap (-n,n)\bigr) \cap (-(n-1), n-1) = A \cap (-(n-1), n-1), \end{equation*}

and

\begin{equation*} \bigl(A \cap (-n,n)\bigr) \setminus (-(n-1), n-1) = A \cap \{x : n-1 \le |x| < n\} = C_n . \end{equation*}

Hence Exercise 8, applied to the set \(A \cap (-n,n)\), gives

\begin{equation*} |A \cap (-n,n)| = |A \cap (-(n-1), n-1)| + |C_n| \qquad \text{for all } n \ge 2 . \end{equation*}

Since \(|A \cap (-1,1)| = |C_1|\), induction on \(n\) now yields

\begin{equation*} |A \cap (-n,n)| = \sum_{k=1}^{n} |C_k| \qquad \text{for all } n \in \mathbb{Z}^{+}, \end{equation*}

and therefore

\begin{equation*} L = \lim_{n \to \infty} |A \cap (-n,n)| = \sum_{k=1}^{\infty} |C_k| . \end{equation*}

Every \(x \in A\) satisfies \(|x| < n\) for some \(n \in \mathbb{Z}^{+}\), and hence lies in \(C_m\) where \(m\) is the least such \(n\); thus \(A = \bigcup_{k=1}^{\infty} C_k\). By the countable subadditivity of outer measure (2.8),

\begin{equation*} |A| \le \sum_{k=1}^{\infty} |C_k| = L . \end{equation*}

Combining the two inequalities gives \(|A| = L = \lim_{t \to \infty} |A \cap (-t,t)|\).

Problem (2A.10)

Prove that \(|[0,1] \setminus \mathbb{Q}| = 1\).

Solution

Write \(E = [0,1] \setminus \mathbb{Q}\), the set of irrational numbers in \([0,1]\).

Upper bound. Since \(E \subseteq [0,1]\), the order-preserving property of outer measure (2.5) together with 2.14 gives

\begin{equation*} |E| \le |[0,1]| = 1 - 0 = 1 . \end{equation*}

Lower bound. The set \(\mathbb{Q}\) is countable, so \([0,1] \cap \mathbb{Q}\) is countable and therefore has outer measure \(0\) by 2.4. Since

\begin{equation*} [0,1] = E \cup \bigl([0,1] \cap \mathbb{Q}\bigr), \end{equation*}

the countable subadditivity of outer measure (2.8), applied to the sequence \(E,\ [0,1] \cap \mathbb{Q},\ \emptyset,\ \emptyset,\dots\), gives

\begin{equation*} 1 = |[0,1]| \le |E| + |[0,1] \cap \mathbb{Q}| = |E| + 0 = |E| , \end{equation*}

where again \(|[0,1]| = 1\) by 2.14.

The two inequalities give \(|[0,1] \setminus \mathbb{Q}| = 1\).

Problem (2A.11)

Prove that if \(I_1, I_2, \dots\) is a disjoint sequence of open intervals, then

\begin{equation*} \Bigl| \bigcup_{k=1}^{\infty} I_k \Bigr| = \sum_{k=1}^{\infty} \ell(I_k). \end{equation*}

Solution

Let \(A = \bigcup_{k=1}^{\infty} I_k\).

Step 0: every open interval \(I\) satisfies \(|I| = \ell(I)\). If \(I = \emptyset\) both sides are \(0\). If \(I = (a,b)\) with \(a,b \in \mathbb{R}\) and \(a < b\), then on the one hand \((a,b) \subseteq [a,b]\), so \(|(a,b)| \le |[a,b]| = b - a\) by 2.5 and 2.14; on the other hand, for every \(\delta\) with \(0 < \delta < (b-a)/2\) we have \([a+\delta, b-\delta] \subseteq (a,b)\), so \(b - a - 2\delta = |[a+\delta,b-\delta]| \le |(a,b)|\) by 2.14 and 2.5, and letting \(\delta \downarrow 0\) gives \(b - a \le |(a,b)|\). Hence \(|(a,b)| = b - a = \ell\bigl((a,b)\bigr)\). If \(I\) is unbounded, then for every \(M > 0\) the interval \(I\) contains a closed interval \([a,b]\) with \(b - a = M\), so \(|I| \ge |[a,b]| = M\) by 2.5 and 2.14; hence \(|I| = \infty = \ell(I)\).

Step 1: the inequality \(|A| \le \sum_{k=1}^{\infty} \ell(I_k)\). The sequence \(I_1, I_2, \dots\) is itself a sequence of open intervals whose union contains \(A\), so the definition of outer measure (2.2) gives \(|A| \le \sum_{k=1}^{\infty} \ell(I_k)\). Disjointness is not needed here.

Step 2: a one-point splitting lemma. I claim that for every \(S \subseteq \mathbb{R}\) and every \(c \in \mathbb{R}\),

\begin{equation*} |S| = \bigl|S \cap (-\infty,c]\bigr| + \bigl|S \cap (c,\infty)\bigr| . \end{equation*}

The inequality \(\le\) follows from countable subadditivity (2.8), since \(S\) is the union of the two sets on the right. For the inequality \(\ge\) we may assume \(|S| < \infty\). Let \(\varepsilon > 0\) and choose open intervals \(J_1, J_2, \dots\) with \(S \subseteq \bigcup_k J_k\) and \(\sum_k \ell(J_k) \le |S| + \varepsilon\). Put \(J_k^{\prime} = J_k \cap (-\infty,c)\) and \(J_k^{\prime\prime} = J_k \cap (c,\infty)\); these are open intervals, and

\begin{equation*} \ell(J_k^{\prime}) + \ell(J_k^{\prime\prime}) = \ell(J_k) \end{equation*}

for each \(k\): this is trivial if \(J_k = \emptyset\); if \(J_k = (a,b)\) is bounded then for \(c \le a\) the pieces are \(\emptyset\) and \((a,b)\), for \(c \ge b\) they are \((a,b)\) and \(\emptyset\), and for \(a < c < b\) they are \((a,c)\) and \((c,b)\) with \((c-a)+(b-c) = b-a\); and if \(J_k\) is unbounded then at least one piece is unbounded, so both sides are \(\infty\). Now

\begin{equation*} S \cap (-\infty,c] \subseteq \Bigl( \bigcup_{k=1}^{\infty} J_k^{\prime} \Bigr) \cup \{c\} \qquad\text{and}\qquad S \cap (c,\infty) \subseteq \bigcup_{k=1}^{\infty} J_k^{\prime\prime} , \end{equation*}

so using 2.8, then 2.3 (a one-point set has outer measure \(0\)) and 2.2,

\begin{equation*} \bigl|S \cap (-\infty,c]\bigr| \le \sum_{k=1}^{\infty} \ell(J_k^{\prime}) \qquad\text{and}\qquad \bigl|S \cap (c,\infty)\bigr| \le \sum_{k=1}^{\infty} \ell(J_k^{\prime\prime}) . \end{equation*}

Adding gives \(|S \cap (-\infty,c]| + |S \cap (c,\infty)| \le \sum_k \ell(J_k) \le |S| + \varepsilon\); since \(\varepsilon > 0\) is arbitrary, the claim follows.

Step 3: splitting at finitely many points. If \(m \in \mathbb{Z}^{+}\), \(c_1 < c_2 < \cdots < c_m\) are real numbers, and \(S \subseteq \mathbb{R}\), then

\begin{equation*} |S| = \bigl|S \cap (-\infty,c_1]\bigr| + \sum_{j=1}^{m-1} \bigl|S \cap (c_j, c_{j+1}]\bigr| + \bigl|S \cap (c_m,\infty)\bigr| . \end{equation*}

Induct on \(m\). The case \(m = 1\) is Step 2. For \(m > 1\), Step 2 at \(c_1\) gives \(|S| = |S \cap (-\infty,c_1]| + |S^{\prime}|\) where \(S^{\prime} = S \cap (c_1,\infty)\). Applying the induction hypothesis to \(S^{\prime}\) with the cut points \(c_2 < \cdots < c_m\) and noting that \(S^{\prime} \cap (-\infty,c_2] = S \cap (c_1,c_2]\), that \(S^{\prime} \cap (c_j,c_{j+1}] = S \cap (c_j,c_{j+1}]\) for \(j \ge 2\), and that \(S^{\prime} \cap (c_m,\infty) = S \cap (c_m,\infty)\), completes the induction.

Step 4: the inequality \(|A| \ge \sum_{k=1}^{n} \ell(I_k)\) for each \(n\). Fix \(n \in \mathbb{Z}^{+}\). If some \(I_k\) with \(k \le n\) is unbounded, then \(|A| \ge |I_k| = \infty\) by 2.5 and Step 0, and the inequality is trivial. So assume \(I_1,\dots,I_n\) are all bounded. Empty intervals contribute \(0\) to \(\sum_{k=1}^{n} \ell(I_k)\), so we may discard them; let \((a_1,b_1),\dots,(a_p,b_p)\) be the nonempty intervals among \(I_1,\dots,I_n\), where \(a_j < b_j\). If \(p = 0\) then \(\sum_{k=1}^{n} \ell(I_k) = 0\) and the inequality is trivial, so assume \(p \ge 1\). If \((a,b)\) and \((a^{\prime},b^{\prime})\) are disjoint nonempty open intervals, then \(b \le a^{\prime}\) or \(b^{\prime} \le a\) (otherwise \(a^{\prime} < b\) and \(a < b^{\prime}\), so \(\max\{a,a^{\prime}\} < \min\{b,b^{\prime}\}\) and the intersection is nonempty). Hence these \(p\) intervals are linearly ordered, and after relabelling

\begin{equation*} a_1 < b_1 \le a_2 < b_2 \le \cdots \le a_p < b_p . \end{equation*}

Let \(c_1 < c_2 < \cdots < c_m\) be the distinct elements of \(\{a_1,b_1,\dots,a_p,b_p\}\), listed in increasing order. For each \(j\), no element of that set lies strictly between \(a_j\) and \(b_j\) (every \(a_i\) and \(b_i\) with \(i < j\) is at most \(a_j\), and every \(a_i\) and \(b_i\) with \(i > j\) is at least \(b_j\)), so \(a_j\) and \(b_j\) are consecutive entries of the list \(c_1,\dots,c_m\). Therefore each interval \((a_j,b_j]\) occurs among the intervals \((c_l, c_{l+1}]\) appearing in Step 3.

Apply Step 3 with \(S = A\) and these cut points, and discard all terms except the \(p\) terms corresponding to the intervals \((a_1,b_1],\dots,(a_p,b_p]\) (all terms are nonnegative, and the discarded terms are distinct from the kept ones because the intervals \((a_j,b_j]\) are distinct members of the partition). This gives

\begin{equation*} |A| \ge \sum_{j=1}^{p} \bigl|A \cap (a_j,b_j]\bigr| . \end{equation*}

Since \((a_j,b_j) \subseteq A \cap (a_j,b_j]\), we get \(|A \cap (a_j,b_j]| \ge |(a_j,b_j)| = b_j - a_j\) by 2.5 and Step 0. Hence

\begin{equation*} |A| \ge \sum_{j=1}^{p} (b_j - a_j) = \sum_{k=1}^{n} \ell(I_k) . \end{equation*}

Step 5: conclusion. Letting \(n \to \infty\) in Step 4 gives \(|A| \ge \sum_{k=1}^{\infty} \ell(I_k)\). Together with Step 1,

\begin{equation*} \Bigl| \bigcup_{k=1}^{\infty} I_k \Bigr| = \sum_{k=1}^{\infty} \ell(I_k), \end{equation*}

as desired.

Problem (2A.12)

Suppose \(r_1, r_2, \dots\) is a sequence that contains every rational number. Let

\begin{equation*} F = \mathbb{R} \setminus \bigcup_{k=1}^{\infty} \Bigl( r_k - \frac{1}{2^k},\ r_k + \frac{1}{2^k} \Bigr). \end{equation*}

(a) Show that \(F\) is a closed subset of \(\mathbb{R}\).

(b) Prove that if \(I\) is an interval contained in \(F\), then \(I\) contains at most one element.

(c) Prove that \(|F| = \infty\).

Solution

Throughout, write

\begin{equation*} G = \bigcup_{k=1}^{\infty} \Bigl( r_k - \frac{1}{2^k},\ r_k + \frac{1}{2^k} \Bigr), \qquad\text{so that}\qquad F = \mathbb{R} \setminus G . \end{equation*}

(a) Each set \(\bigl( r_k - 2^{-k}, r_k + 2^{-k} \bigr)\) is an open subset of \(\mathbb{R}\), and an arbitrary union of open sets is open; hence \(G\) is open. Therefore \(F = \mathbb{R} \setminus G\) is the complement of an open set, which is precisely what it means for \(F\) to be closed.

(b) Suppose, for contradiction, that \(I \subseteq F\) is an interval containing two distinct elements \(x\) and \(y\), say \(x < y\). Because \(I\) is an interval containing \(x\) and \(y\), we have \([x,y] \subseteq I\). By the density of the rationals in \(\mathbb{R}\) there is a rational number \(q\) with \(x < q < y\), so \(q \in I \subseteq F\).

On the other hand, the sequence \(r_1, r_2, \dots\) contains every rational number, so \(q = r_j\) for some \(j \in \mathbb{Z}^{+}\). Then

\begin{equation*} q = r_j \in \Bigl( r_j - \frac{1}{2^{\,j}},\ r_j + \frac{1}{2^{\,j}} \Bigr) \subseteq G , \end{equation*}

so \(q \notin F\). This contradiction shows that \(I\) contains at most one element.

(c) First bound the outer measure of \(G\). The sequence of open intervals \(\bigl( r_k - 2^{-k}, r_k + 2^{-k} \bigr)\), \(k \in \mathbb{Z}^{+}\), has union \(G\), so by the definition of outer measure (2.2),

\begin{equation*} |G| \le \sum_{k=1}^{\infty} \ell\Bigl( \Bigl( r_k - \frac{1}{2^k},\ r_k + \frac{1}{2^k} \Bigr) \Bigr) = \sum_{k=1}^{\infty} \frac{2}{2^k} = 2 . \end{equation*}

Now fix \(n \in \mathbb{Z}^{+}\). Since \(F \cup G = \mathbb{R}\), we have \([-n,n] \subseteq F \cup G\), and hence by 2.5 and the subadditivity of outer measure (2.8),

\begin{equation*} |[-n,n]| \le |F \cup G| \le |F| + |G| \le |F| + 2 , \end{equation*}

the first inequality by 2.5, the second by 2.8 applied to the sequence \(F, G, \emptyset, \emptyset, \dots\). By 2.14 we have \(|[-n,n]| = 2n\), and since \(|G| \le 2 < \infty\) we may subtract, obtaining

\begin{equation*} |F| \ge 2n - 2 . \end{equation*}

This holds for every \(n \in \mathbb{Z}^{+}\), so \(|F| = \infty\).

It is worth pausing over what (b) and (c) say together: \(F\) is a closed set that contains no interval of positive length whatsoever (so \(F\) has empty interior and is nowhere dense), and yet \(F\) has infinite outer measure. Outer measure and topological size are genuinely different notions.

Problem (2A.13)

Suppose \(\varepsilon > 0\). Prove that there exists a subset \(F\) of \([0,1]\) such that \(F\) is closed, every element of \(F\) is an irrational number, and \(|F| > 1 - \varepsilon\).

Solution

The set \([0,1] \cap \mathbb{Q}\) is countably infinite, so we may list it as a sequence \(q_1, q_2, \dots\) containing every rational number in \([0,1]\). Define

\begin{equation*} G = \bigcup_{k=1}^{\infty} \Bigl( q_k - \frac{\varepsilon}{2^{k+2}},\ q_k + \frac{\varepsilon}{2^{k+2}} \Bigr) \qquad\text{and}\qquad F = [0,1] \setminus G . \end{equation*}

\(F\) is closed. Each interval in the union defining \(G\) is open, so \(G\) is open, so \(\mathbb{R} \setminus G\) is closed. Since \([0,1]\) is closed and \(F = [0,1] \cap (\mathbb{R} \setminus G)\) is the intersection of two closed sets, \(F\) is closed. Also \(F \subseteq [0,1]\) by construction.

Every element of \(F\) is irrational. Suppose \(x \in F\) were rational. Then \(x \in [0,1] \cap \mathbb{Q}\), so \(x = q_k\) for some \(k \in \mathbb{Z}^{+}\), and hence

\begin{equation*} x = q_k \in \Bigl( q_k - \frac{\varepsilon}{2^{k+2}},\ q_k + \frac{\varepsilon}{2^{k+2}} \Bigr) \subseteq G , \end{equation*}

contradicting \(x \notin G\). Thus every element of \(F\) is irrational.

The estimate \(|F| > 1 - \varepsilon\). The intervals used to define \(G\) form a sequence of open intervals whose union is \(G\), so by the definition of outer measure (2.2),

\begin{equation*} |G| \le \sum_{k=1}^{\infty} \frac{2\varepsilon}{2^{k+2}} = \frac{\varepsilon}{2} \sum_{k=1}^{\infty} \frac{1}{2^{k}} = \frac{\varepsilon}{2} . \end{equation*}

Since \([0,1] \subseteq F \cup G\) (a point of \([0,1]\) either lies in \(G\) or, failing that, lies in \(F\)), the order-preserving property (2.5) and the subadditivity of outer measure (2.8), together with \(|[0,1]| = 1\) from 2.14, give

\begin{equation*} 1 = |[0,1]| \le |F \cup G| \le |F| + |G| \le |F| + \frac{\varepsilon}{2} . \end{equation*}

Because \(|G| \le \varepsilon/2 < \infty\), we may subtract it. Hence

\begin{equation*} |F| \ge 1 - \frac{\varepsilon}{2} > 1 - \varepsilon , \end{equation*}

the strict inequality because \(\varepsilon > 0\). This \(F\) has all three required properties.

Note that \(F\) contains no interval with two distinct elements: any such interval would contain a rational number, by the density of \(\mathbb{Q}\), and \(F\) contains no rational number. So \(F\) is a closed, nowhere dense subset of \([0,1]\) whose outer measure is nevertheless as close to \(1\) as we please. Contrast this with the fact that a countable set, however dense, has outer measure \(0\) by 2.4.

Problem (2A.14)

Consider the following figure, which is drawn accurately to scale.

[Description of the figure. In the coordinate plane, with the horizontal axis marked at \(2, 5, 8, 11, 14, 17, 20\) and the vertical axis marked at \(1, 3, 5, 7, 9\), three colored pieces are drawn: a yellow right triangle with vertices \((0,0)\), \((11,0)\), \((11,5)\); a red rectangle with vertices \((11,0)\), \((20,0)\), \((20,5)\), \((11,5)\); and a blue right triangle with vertices \((11,5)\), \((20,5)\), \((20,9)\). Their upper boundary runs from \((0,0)\) to \((11,5)\) to \((20,9)\) and looks to the eye like a single straight segment, so the three pieces together appear to fill the right triangle with vertices \((0,0)\), \((20,0)\), \((20,9)\).]

(a) Show that the right triangle whose vertices are \((0,0)\), \((20,0)\), and \((20,9)\) has area \(90\).

[We have not defined area yet, but just use the elementary formulas for the areas of triangles and rectangles that you learned long ago.]

(b) Show that the yellow (lower) right triangle has area \(27.5\).

(c) Show that the red rectangle has area \(45\).

(d) Show that the blue (upper) right triangle has area \(18\).

(e) Add the results of parts (b), (c), and (d), showing that the area of the colored region is \(90.5\).

(f) Seeing the figure above, most people expect parts (a) and (e) to have the same result. Yet in (a) we found area \(90\), and in (e) we found area \(90.5\). Explain why these results differ.

[You may be tempted to think that what we have here is a two-dimensional example similar to the result about the nonadditivity of outer measure (2.18). However, genuine examples of nonadditivity require much more complicated sets than in this example.]

Solution

(a) Call this triangle \(T\). The side from \((0,0)\) to \((20,0)\) is horizontal and the side from \((20,0)\) to \((20,9)\) is vertical, so \(T\) has its right angle at \((20,0)\), with legs of length \(20\) and \(9\). Hence

\begin{equation*} \operatorname{area}(T) = \tfrac{1}{2} \cdot 20 \cdot 9 = 90 . \end{equation*}

(b) The yellow triangle has vertices \((0,0)\), \((11,0)\), \((11,5)\), with its right angle at \((11,0)\) and legs of length \(11\) and \(5\). Hence its area is

\begin{equation*} \tfrac{1}{2} \cdot 11 \cdot 5 = \tfrac{55}{2} = 27.5 . \end{equation*}

(c) The red rectangle has vertices \((11,0)\), \((20,0)\), \((20,5)\), \((11,5)\); its horizontal side has length \(20 - 11 = 9\) and its vertical side has length \(5 - 0 = 5\). Hence its area is

\begin{equation*} 9 \cdot 5 = 45 . \end{equation*}

(d) The blue triangle has vertices \((11,5)\), \((20,5)\), \((20,9)\), with its right angle at \((20,5)\) and legs of length \(20 - 11 = 9\) and \(9 - 5 = 4\). Hence its area is

\begin{equation*} \tfrac{1}{2} \cdot 9 \cdot 4 = 18 . \end{equation*}

(e) The three colored pieces overlap only along parts of their boundaries: the yellow triangle meets the red rectangle in the segment \(\{11\} \times [0,5]\), the red rectangle meets the blue triangle in the segment \([11,20] \times \{5\}\), and the yellow and blue triangles meet only at the single point \((11,5)\). Segments and points contribute no area, so the area of the colored region is the sum of the three areas:

\begin{equation*} 27.5 + 45 + 18 = 90.5 . \end{equation*}

(f) The colored region is simply not the triangle of part (a). There is no paradox, and no failure of additivity; the mistake is the visual assumption that the three pieces fit together to fill \(T\).

The union \(Q\) of the three colored pieces is the quadrilateral with vertices \((0,0)\), \((20,0)\), \((20,9)\), \((11,5)\), taken in that (cyclic) order. The point \((11,5)\) is a genuine corner of \(Q\), not a point on a straight edge. Indeed, the line through \((0,0)\) and \((20,9)\) is

\begin{equation*} y = \tfrac{9}{20} x , \end{equation*}

and at \(x = 11\) this line has height \(\tfrac{99}{20} = 4.95\), whereas the colored region reaches up to height \(5\) there. Equivalently, compare slopes: the hypotenuse of the yellow triangle has slope \(\tfrac{5}{11} = 0.4545\ldots\), the hypotenuse of the blue triangle has slope \(\tfrac{4}{9} = 0.4444\ldots\), and the hypotenuse of \(T\) has slope \(\tfrac{9}{20} = 0.45\). These three numbers are distinct, so the upper boundary of the colored region is a bent path, not a line segment. The bend is tiny (the two slopes differ by about \(0.01\)), which is exactly why the figure is deceptive even though it is drawn accurately to scale.

Since \(5 > 4.95\), the bend goes upward, so \(T \subseteq Q\) and \(Q\) is strictly larger than \(T\). The difference \(Q \setminus T\) is the very thin triangle \(S\) with vertices \((0,0)\), \((11,5)\), \((20,9)\). Its area is computed as follows: take as base the segment from \((0,0)\) to \((20,9)\), of length \(\sqrt{20^2 + 9^2} = \sqrt{481}\); the corresponding height is the distance from \((11,5)\) to the line \(9x - 20y = 0\), namely

\begin{equation*} \frac{|9 \cdot 11 - 20 \cdot 5|}{\sqrt{9^2 + 20^2}} = \frac{1}{\sqrt{481}} . \end{equation*}

Hence

\begin{equation*} \operatorname{area}(S) = \tfrac{1}{2} \cdot \sqrt{481} \cdot \frac{1}{\sqrt{481}} = \tfrac{1}{2} , \end{equation*}

and \(\operatorname{area}(Q) = \operatorname{area}(T) + \operatorname{area}(S) = 90 + 0.5 = 90.5\), in agreement with part (e). (The same value follows from the shoelace formula applied to the vertices \((0,0)\), \((20,0)\), \((20,9)\), \((11,5)\).)

Finally, to address the bracketed remark: this example has nothing to do with the nonadditivity of outer measure (2.18). Here the three pieces are nonoverlapping and their areas do add up to the area of their union, which is \(90.5\). What fails is only the false identification of that union with the triangle \(T\). The genuinely nonadditive sets of 2.18 are produced by the Axiom of Choice and are far too pathological to be drawn; no picture assembled from polygons could ever exhibit nonadditivity.

Exercises 2B

Problem (2B.1)

Show that \(\mathcal{S} = \left\{\bigcup_{n \in K} (n, n+1] : K \subseteq \mathbb{Z}\right\}\) is a \(\sigma\)-algebra on \(\mathbb{R}\).

Solution

For \(K \subseteq \mathbb{Z}\) write

\begin{equation*} E_K = \bigcup_{n \in K} (n, n+1], \end{equation*}

so that \(\mathcal{S} = \{E_K : K \subseteq \mathbb{Z}\}\). Everything follows from the observation that the intervals \((n, n+1]\), for \(n \in \mathbb{Z}\), form a partition of \(\mathbb{R}\).

Step 1: each \(x \in \mathbb{R}\) lies in exactly one of the intervals \((n, n+1]\).

Existence: let \(n_x = \lceil x \rceil - 1\), where \(\lceil x \rceil\) denotes the smallest integer that is greater than or equal to \(x\). If \(x \in \mathbb{Z}\) then \(\lceil x \rceil = x\) and \(n_x = x - 1\), so \(n_x < x \le n_x + 1\). If \(x \notin \mathbb{Z}\) then \(\lceil x \rceil - 1 < x < \lceil x \rceil\), so again \(n_x < x \le n_x + 1\). Either way \(x \in (n_x, n_x + 1]\).

Uniqueness: suppose \(m, n \in \mathbb{Z}\), \(m < n\), and \(x \in (m, m+1] \cap (n, n+1]\). Since \(m < n\) and both are integers, \(n \ge m + 1\). Thus \(x \le m + 1 \le n\), contradicting \(x > n\). Hence no \(x\) lies in two of these intervals, and \(n_x\) is the unique integer with \(x \in (n_x, n_x+1]\).

Step 2: for every \(K \subseteq \mathbb{Z}\) and every \(x \in \mathbb{R}\),

\begin{equation*} x \in E_K \iff n_x \in K . \end{equation*}

Indeed, if \(n_x \in K\) then \(x \in (n_x, n_x+1] \subseteq E_K\). Conversely, if \(x \in E_K\) then \(x \in (n, n+1]\) for some \(n \in K\), and Step 1 forces \(n = n_x\), so \(n_x \in K\). In particular \(K \mapsto E_K\) is injective, though we shall not need that.

Now verify the three conditions in the definition of a \(\sigma\)-algebra (2.23).

The empty set: taking \(K = \emptyset\) gives \(E_\emptyset = \emptyset\) (an empty union), so \(\emptyset \in \mathcal{S}\).

Closure under complementation: we claim that

\begin{equation*} \mathbb{R} \setminus E_K = E_{\mathbb{Z} \setminus K} \end{equation*}

for every \(K \subseteq \mathbb{Z}\). Let \(x \in \mathbb{R}\). By Step 2, \(x \notin E_K\) if and only if \(n_x \notin K\), which (as \(n_x \in \mathbb{Z}\)) holds if and only if \(n_x \in \mathbb{Z} \setminus K\), which by Step 2 again holds if and only if \(x \in E_{\mathbb{Z} \setminus K}\). This proves the claim, and since \(\mathbb{Z} \setminus K \subseteq \mathbb{Z}\) we conclude \(\mathbb{R} \setminus E_K \in \mathcal{S}\).

Closure under countable unions: suppose \(E_{K_1}, E_{K_2}, \ldots\) is a sequence of elements of \(\mathcal{S}\), with each \(K_j \subseteq \mathbb{Z}\). Put \(K = \bigcup_{j=1}^\infty K_j \subseteq \mathbb{Z}\). Then for \(x \in \mathbb{R}\), using Step 2 three times,

\begin{equation*} x \in \bigcup_{j=1}^\infty E_{K_j} \iff n_x \in K_j \text{ for some } j \iff n_x \in K \iff x \in E_K . \end{equation*}

Hence \(\bigcup_{j=1}^\infty E_{K_j} = E_K \in \mathcal{S}\).

Thus \(\mathcal{S}\) contains \(\emptyset\), is closed under complementation in \(\mathbb{R}\), and is closed under countable unions; that is, \(\mathcal{S}\) is a \(\sigma\)-algebra on \(\mathbb{R}\).

(Note in passing that \(\mathbb{R} = E_{\mathbb{Z}} \in \mathcal{S}\), in agreement with 2.25(a), and that every element of \(\mathcal{S}\) is a Borel set, since each \((n, n+1]\) is a Borel set and the unions involved are countable. So \(\mathcal{S}\) is a \(\sigma\)-algebra contained in the Borel \(\sigma\)-algebra.)

Problem (2B.2)

Verify both bullet points in Example 2.28, which assert the following.

  • (a) Suppose \(X\) is a set and \(\mathcal{A}\) is the set of subsets of \(X\) that consist of exactly one element, \(\mathcal{A} = \bigl\{\{x\} : x \in X\bigr\}\). Then the smallest \(\sigma\)-algebra on \(X\) containing \(\mathcal{A}\) is the set of all subsets \(E\) of \(X\) such that \(E\) is countable or \(X \setminus E\) is countable.

  • (b) Suppose \(\mathcal{A} = \{(0,1), (0,\infty)\}\). Then the smallest \(\sigma\)-algebra on \(\mathbb{R}\) containing \(\mathcal{A}\) is \(\{\emptyset,\ (0,1),\ (0,\infty),\ (-\infty,0] \cup [1,\infty),\ (-\infty,0],\ [1,\infty),\ (-\infty,1),\ \mathbb{R}\}\).

Solution

Throughout we use the standard strategy for identifying a smallest \(\sigma\)-algebra containing a collection \(\mathcal{A}\): exhibit a candidate collection \(\mathcal{T}\), show (i) \(\mathcal{T}\) is a \(\sigma\)-algebra, (ii) \(\mathcal{A} \subseteq \mathcal{T}\), and (iii) every \(\sigma\)-algebra containing \(\mathcal{A}\) contains \(\mathcal{T}\). By 2.27 the smallest \(\sigma\)-algebra containing \(\mathcal{A}\) is the intersection of all \(\sigma\)-algebras containing \(\mathcal{A}\); (i) and (ii) say \(\mathcal{T}\) is one of the collections being intersected, so the intersection is contained in \(\mathcal{T}\), while (iii) says \(\mathcal{T}\) is contained in each of them, hence in the intersection. Therefore \(\mathcal{T}\) is exactly the smallest such \(\sigma\)-algebra.

Part (a). Let

\begin{equation*} \mathcal{T} = \{E \subseteq X : E \text{ is countable or } X \setminus E \text{ is countable}\}. \end{equation*}

(i) \(\mathcal{T}\) is a \(\sigma\)-algebra on \(X\). First, \(\emptyset\) is countable, so \(\emptyset \in \mathcal{T}\). Next, the defining condition on \(E\) is symmetric in \(E\) and \(X \setminus E\) (because \(X \setminus (X \setminus E) = E\)), so \(E \in \mathcal{T}\) implies \(X \setminus E \in \mathcal{T}\). Finally, suppose \(E_1, E_2, \ldots \in \mathcal{T}\) and let \(E = \bigcup_{k=1}^\infty E_k\). If every \(E_k\) is countable, then \(E\) is a countable union of countable sets and hence is countable, so \(E \in \mathcal{T}\). Otherwise there is some \(m\) with \(E_m\) uncountable; since \(E_m \in \mathcal{T}\), the set \(X \setminus E_m\) must be countable. Because \(E \supseteq E_m\) we have \(X \setminus E \subseteq X \setminus E_m\), so \(X \setminus E\) is countable and again \(E \in \mathcal{T}\). Thus \(\mathcal{T}\) is a \(\sigma\)-algebra.

(ii) \(\mathcal{A} \subseteq \mathcal{T}\), since each \(\{x\}\) is finite and therefore countable.

(iii) Suppose \(\mathcal{S}\) is a \(\sigma\)-algebra on \(X\) with \(\mathcal{A} \subseteq \mathcal{S}\), and suppose \(E \in \mathcal{T}\). If \(E\) is countable, then either \(E = \emptyset \in \mathcal{S}\), or we may list \(E = \{x_1, x_2, \ldots\}\) (a finite or infinite list) and write

\begin{equation*} E = \bigcup_{k} \{x_k\}, \end{equation*}

a countable union of elements of \(\mathcal{A} \subseteq \mathcal{S}\); hence \(E \in \mathcal{S}\) (pad a finite list with copies of \(\emptyset\) to get a sequence). If instead \(X \setminus E\) is countable, the case just treated gives \(X \setminus E \in \mathcal{S}\), and then \(E = X \setminus (X \setminus E) \in \mathcal{S}\). So \(\mathcal{T} \subseteq \mathcal{S}\).

By the strategy above, \(\mathcal{T}\) is the smallest \(\sigma\)-algebra on \(X\) containing \(\mathcal{A}\), which is the first bullet point.

Part (b). Set

\begin{equation*} A_1 = (0,1), \qquad A_2 = [1,\infty), \qquad A_3 = (-\infty, 0]. \end{equation*}

These three sets are nonempty, pairwise disjoint, and their union is \(\mathbb{R}\); that is, \(\{A_1, A_2, A_3\}\) is a partition of \(\mathbb{R}\). For \(K \subseteq \{1,2,3\}\) put \(E_K = \bigcup_{i \in K} A_i\). Listing the eight subsets \(K\) gives

\begin{equation*} \begin{aligned} E_\emptyset &= \emptyset, & E_{\{1\}} &= (0,1), & E_{\{2\}} &= [1,\infty), & E_{\{3\}} &= (-\infty,0],\\ E_{\{1,2\}} &= (0,\infty), & E_{\{1,3\}} &= (-\infty,1), & E_{\{2,3\}} &= (-\infty,0] \cup [1,\infty), & E_{\{1,2,3\}} &= \mathbb{R}, \end{aligned} \end{equation*}

so the collection \(\mathcal{T} = \{E_K : K \subseteq \{1,2,3\}\}\) is precisely the eight-element collection displayed in the exercise. (The eight sets are distinct: since the \(A_i\) are nonempty and pairwise disjoint, \(E_K\) determines \(K\) via \(i \in K \iff A_i \subseteq E_K\).)

(i) \(\mathcal{T}\) is a \(\sigma\)-algebra on \(\mathbb{R}\). Since \(\{A_1,A_2,A_3\}\) partitions \(\mathbb{R}\), each \(x \in \mathbb{R}\) lies in exactly one \(A_{i(x)}\), and \(x \in E_K \iff i(x) \in K\). Exactly as in Exercise 1 this yields

\begin{equation*} \mathbb{R} \setminus E_K = E_{\{1,2,3\} \setminus K} \qquad\text{and}\qquad \bigcup_{j=1}^\infty E_{K_j} = E_{K_1 \cup K_2 \cup \cdots}, \end{equation*}

and \(\emptyset = E_\emptyset \in \mathcal{T}\). Hence \(\mathcal{T}\) is closed under complementation and countable unions and contains \(\emptyset\), so it is a \(\sigma\)-algebra.

(ii) \(\mathcal{A} \subseteq \mathcal{T}\), because \((0,1) = E_{\{1\}}\) and \((0,\infty) = E_{\{1,2\}}\).

(iii) Suppose \(\mathcal{S}\) is a \(\sigma\)-algebra on \(\mathbb{R}\) with \((0,1), (0,\infty) \in \mathcal{S}\). Then

\begin{equation*} A_1 = (0,1) \in \mathcal{S}, \qquad A_2 = [1,\infty) = (0,\infty) \setminus (0,1) \in \mathcal{S}, \qquad A_3 = (-\infty,0] = \mathbb{R} \setminus (0,\infty) \in \mathcal{S}, \end{equation*}

using 2.25(a) and 2.25(b). Since \(\mathcal{S}\) is closed under finite unions (2.25(b)), every \(E_K\) is in \(\mathcal{S}\). Thus \(\mathcal{T} \subseteq \mathcal{S}\).

By the strategy above, \(\mathcal{T}\) is the smallest \(\sigma\)-algebra on \(\mathbb{R}\) containing \(\{(0,1), (0,\infty)\}\), which is the second bullet point.

Problem (2B.3)

Suppose \(\mathcal{S}\) is the smallest \(\sigma\)-algebra on \(\mathbb{R}\) containing \(\{(r,s] : r, s \in \mathbb{Q}\}\). Prove that \(\mathcal{S}\) is the collection of Borel subsets of \(\mathbb{R}\).

Solution

Write \(\mathcal{B}\) for the collection of Borel subsets of \(\mathbb{R}\), which by 2.29 is the smallest \(\sigma\)-algebra on \(\mathbb{R}\) containing all open subsets of \(\mathbb{R}\). We prove the two inclusions \(\mathcal{S} \subseteq \mathcal{B}\) and \(\mathcal{B} \subseteq \mathcal{S}\).

Step 1: \(\mathcal{S} \subseteq \mathcal{B}\).

It suffices to show that \((r,s] \in \mathcal{B}\) for all \(r, s \in \mathbb{Q}\), for then \(\mathcal{B}\) is a \(\sigma\)-algebra on \(\mathbb{R}\) containing \(\{(r,s] : r, s \in \mathbb{Q}\}\), and \(\mathcal{S}\), being the smallest such \(\sigma\)-algebra, is contained in \(\mathcal{B}\) (see the remark following 2.27).

If \(r \ge s\) then \((r,s] = \emptyset \in \mathcal{B}\). If \(r < s\) then

\begin{equation*} (r,s] = \bigcap_{k=1}^\infty \left(r,\ s + \tfrac{1}{k}\right), \end{equation*}

because a real number \(x\) satisfies \(r < x \le s\) if and only if \(r < x < s + \frac1k\) for every \(k \in \mathbb{Z}^+\) (the forward implication is clear; conversely, \(x < s + \frac1k\) for all \(k\) forces \(x \le s\)). Each set \(\left(r, s + \frac1k\right)\) is open, hence lies in \(\mathcal{B}\), and \(\mathcal{B}\) is closed under countable intersections by 2.25(c). Thus \((r,s] \in \mathcal{B}\).

Step 2: \(\mathcal{B} \subseteq \mathcal{S}\).

Since \(\mathcal{B}\) is the smallest \(\sigma\)-algebra containing the open subsets of \(\mathbb{R}\), it suffices to show that every open subset of \(\mathbb{R}\) belongs to \(\mathcal{S}\).

So let \(G \subseteq \mathbb{R}\) be open, and let

\begin{equation*} P = \{(r,s) \in \mathbb{Q} \times \mathbb{Q} : (r,s] \subseteq G\}. \end{equation*}

We claim

\begin{equation*} G = \bigcup_{(r,s) \in P} (r,s]. \end{equation*}

The inclusion \(\supseteq\) holds because every set in the union is contained in \(G\) by definition of \(P\). For the inclusion \(\subseteq\), suppose \(x \in G\). Because \(G\) is open there exists \(\varepsilon > 0\) with \((x - \varepsilon, x + \varepsilon) \subseteq G\). By the density of \(\mathbb{Q}\) in \(\mathbb{R}\) choose rationals \(r\) and \(s\) with

\begin{equation*} x - \varepsilon < r < x \qquad\text{and}\qquad x < s < x + \varepsilon . \end{equation*}

Then \(x \in (r,s]\) and \((r,s] \subseteq (x-\varepsilon, x+\varepsilon) \subseteq G\), so \((r,s) \in P\) and \(x\) belongs to the right side. This proves the claim.

Now \(P \subseteq \mathbb{Q} \times \mathbb{Q}\) is a countable set, so the union above is a countable union of elements of \(\mathcal{S}\) (if \(P = \emptyset\) then \(G = \emptyset \in \mathcal{S}\); otherwise list \(P\) as a finite or infinite sequence, padding with repetitions if necessary to obtain a sequence). Because \(\mathcal{S}\) is closed under countable unions, \(G \in \mathcal{S}\).

Thus \(\mathcal{S}\) is a \(\sigma\)-algebra on \(\mathbb{R}\) containing every open subset of \(\mathbb{R}\), so \(\mathcal{B} \subseteq \mathcal{S}\).

Combining the two steps, \(\mathcal{S} = \mathcal{B}\); that is, \(\mathcal{S}\) is exactly the collection of Borel subsets of \(\mathbb{R}\).

Problem (2B.4)

Suppose \(\mathcal{S}\) is the smallest \(\sigma\)-algebra on \(\mathbb{R}\) containing \(\{(r,n] : r \in \mathbb{Q},\ n \in \mathbb{Z}\}\). Prove that \(\mathcal{S}\) is the collection of Borel subsets of \(\mathbb{R}\).

Solution

Let \(\mathcal{B}\) denote the collection of Borel subsets of \(\mathbb{R}\).

Step 1: \(\mathcal{S} \subseteq \mathcal{B}\).

Every integer is rational, so \(\{(r,n] : r \in \mathbb{Q},\, n \in \mathbb{Z}\} \subseteq \{(r,s] : r,s \in \mathbb{Q}\}\), and Step 1 of Exercise 3 showed each such half-open interval is a Borel set: either it is empty, or

\begin{equation*} (r,n] = \bigcap_{k=1}^\infty \left(r,\ n + \tfrac1k\right), \end{equation*}

a countable intersection of open sets, which lies in \(\mathcal{B}\) by 2.25(c). Hence \(\mathcal{B}\) is a \(\sigma\)-algebra containing the generating collection, and therefore contains the smallest such \(\sigma\)-algebra, namely \(\mathcal{S}\).

Step 2: \(\mathcal{B} \subseteq \mathcal{S}\).

By Exercise 3, the smallest \(\sigma\)-algebra on \(\mathbb{R}\) containing \(\{(r,s] : r,s \in \mathbb{Q}\}\) equals \(\mathcal{B}\). So it suffices to prove that \((r,s] \in \mathcal{S}\) for all \(r,s \in \mathbb{Q}\); the smallest \(\sigma\)-algebra containing all such intervals is then contained in \(\mathcal{S}\), giving \(\mathcal{B} \subseteq \mathcal{S}\).

Let \(r, s \in \mathbb{Q}\). If \(r \ge s\) then \((r,s] = \emptyset \in \mathcal{S}\). So assume \(r < s\), and choose an integer \(n\) with \(n \ge s\) (possible by the Archimedean property). We claim

\begin{equation*} (r,s] = (r,n] \setminus (s,n]. \end{equation*}

To see this, let \(x \in \mathbb{R}\). If \(r < x \le s\), then, since \(s \le n\), we have \(r < x \le n\), so \(x \in (r,n]\); and \(x \le s\) shows \(x \notin (s,n]\). Conversely, suppose \(x \in (r,n]\) and \(x \notin (s,n]\). From \(x \in (r,n]\) we get \(r < x \le n\); the failure of \(s < x \le n\) then forces \(x \le s\). Hence \(r < x \le s\). This proves the claim.

Both \((r,n]\) and \((s,n]\) are in \(\mathcal{S}\) (they are among the generators, as \(r, s \in \mathbb{Q}\) and \(n \in \mathbb{Z}\)), and \(\sigma\)-algebras are closed under set differences by 2.25(b). Therefore \((r,s] \in \mathcal{S}\).

Combining the two steps gives \(\mathcal{S} = \mathcal{B}\), so \(\mathcal{S}\) is the collection of Borel subsets of \(\mathbb{R}\).

Problem (2B.5)

Suppose \(\mathcal{S}\) is the smallest \(\sigma\)-algebra on \(\mathbb{R}\) containing \(\{(r, r+1) : r \in \mathbb{Q}\}\). Prove that \(\mathcal{S}\) is the collection of Borel subsets of \(\mathbb{R}\).

Solution

Let \(\mathcal{B}\) denote the collection of Borel subsets of \(\mathbb{R}\).

Step 1: \(\mathcal{S} \subseteq \mathcal{B}\).

Each interval \((r, r+1)\) is an open subset of \(\mathbb{R}\) and hence is a Borel set. Thus \(\mathcal{B}\) is a \(\sigma\)-algebra on \(\mathbb{R}\) containing \(\{(r,r+1) : r \in \mathbb{Q}\}\), and therefore contains the smallest such \(\sigma\)-algebra, which is \(\mathcal{S}\).

Step 2: short rational intervals lie in \(\mathcal{S}\).

We claim that if \(p, q \in \mathbb{Q}\) satisfy \(p < q\) and \(q - p \le 1\), then \((p,q) \in \mathcal{S}\). Indeed, both \(p\) and \(q-1\) are rational, so

\begin{equation*} (p, p+1) \in \mathcal{S} \qquad\text{and}\qquad (q-1, q) = \bigl((q-1),\ (q-1)+1\bigr) \in \mathcal{S}, \end{equation*}

and \(\sigma\)-algebras are closed under finite intersections by 2.25(b). Now

\begin{equation*} (p, p+1) \cap (q-1, q) = \bigl(\max\{p,\ q-1\},\ \min\{p+1,\ q\}\bigr) = (p, q), \end{equation*}

where the last equality uses \(q - p \le 1\), which gives \(q - 1 \le p\) and \(q \le p+1\). Hence \((p,q) \in \mathcal{S}\).

Step 3: \(\mathcal{B} \subseteq \mathcal{S}\).

Since \(\mathcal{B}\) is by 2.29 the smallest \(\sigma\)-algebra containing the open subsets of \(\mathbb{R}\), it is enough to show that every open \(G \subseteq \mathbb{R}\) lies in \(\mathcal{S}\).

Let \(G \subseteq \mathbb{R}\) be open and set

\begin{equation*} P = \{(p,q) \in \mathbb{Q} \times \mathbb{Q} : p < q,\ q - p \le 1,\ (p,q) \subseteq G\}. \end{equation*}

We claim \(G = \bigcup_{(p,q) \in P} (p,q)\). The inclusion \(\supseteq\) is immediate from the definition of \(P\). For \(\subseteq\), let \(x \in G\). Because \(G\) is open there exists \(\varepsilon > 0\) with \((x - \varepsilon, x + \varepsilon) \subseteq G\), and we may shrink \(\varepsilon\) so that \(\varepsilon < \frac12\). By density of \(\mathbb{Q}\) in \(\mathbb{R}\) choose rationals \(p, q\) with

\begin{equation*} x - \varepsilon < p < x \qquad\text{and}\qquad x < q < x + \varepsilon . \end{equation*}

Then \(p < x < q\), and \(q - p < 2\varepsilon < 1\), and \((p,q) \subseteq (x-\varepsilon, x+\varepsilon) \subseteq G\). So \((p,q) \in P\) and \(x \in (p,q)\), proving the claim.

The index set \(P\) is a subset of the countable set \(\mathbb{Q} \times \mathbb{Q}\), so the union above is a countable union (if \(P = \emptyset\) then \(G = \emptyset \in \mathcal{S}\); otherwise enumerate \(P\) as a sequence, repeating terms if \(P\) is finite). By Step 2 each set in the union is in \(\mathcal{S}\), and \(\mathcal{S}\) is closed under countable unions. Hence \(G \in \mathcal{S}\).

Therefore \(\mathcal{S}\) is a \(\sigma\)-algebra containing every open subset of \(\mathbb{R}\), so \(\mathcal{B} \subseteq \mathcal{S}\). Combined with Step 1 this gives \(\mathcal{S} = \mathcal{B}\), so \(\mathcal{S}\) is the collection of Borel subsets of \(\mathbb{R}\).

Problem (2B.6)

Suppose \(\mathcal{S}\) is the smallest \(\sigma\)-algebra on \(\mathbb{R}\) containing \(\{[r,\infty) : r \in \mathbb{Q}\}\). Prove that \(\mathcal{S}\) is the collection of Borel subsets of \(\mathbb{R}\).

Solution

Let \(\mathcal{B}\) denote the collection of Borel subsets of \(\mathbb{R}\).

Step 1: \(\mathcal{S} \subseteq \mathcal{B}\).

For each \(r \in \mathbb{Q}\) the set \([r,\infty)\) is a closed subset of \(\mathbb{R}\), because its complement \((-\infty, r)\) is open. By the first bullet point of 2.30, every closed subset of \(\mathbb{R}\) is a Borel set. Hence \(\mathcal{B}\) is a \(\sigma\)-algebra on \(\mathbb{R}\) containing \(\{[r,\infty) : r \in \mathbb{Q}\}\), and therefore contains the smallest such \(\sigma\)-algebra, namely \(\mathcal{S}\).

Step 2: bounded open intervals with rational endpoints lie in \(\mathcal{S}\).

First, for each \(p \in \mathbb{Q}\),

\begin{equation*} (p, \infty) = \bigcup_{k=1}^\infty \left[p + \tfrac1k,\ \infty\right). \end{equation*}

Indeed, if \(x \ge p + \frac1k\) for some \(k \in \mathbb{Z}^+\) then \(x > p\); conversely, if \(x > p\) then by the Archimedean property there is \(k \in \mathbb{Z}^+\) with \(\frac1k \le x - p\), so \(x \in \left[p + \frac1k, \infty\right)\). Each \(p + \frac1k\) is rational, so each set in this union is a generator of \(\mathcal{S}\); as \(\mathcal{S}\) is closed under countable unions, \((p,\infty) \in \mathcal{S}\).

Second, for each \(q \in \mathbb{Q}\),

\begin{equation*} (-\infty, q) = \mathbb{R} \setminus [q, \infty) \in \mathcal{S}, \end{equation*}

since \(\mathcal{S}\) is closed under complementation.

Therefore, for all \(p, q \in \mathbb{Q}\),

\begin{equation*} (p,q) = (p,\infty) \cap (-\infty, q) \in \mathcal{S}, \end{equation*}

using closure of \(\mathcal{S}\) under finite intersections (2.25(b)). (When \(p \ge q\) this set is empty, which is also in \(\mathcal{S}\).)

Step 3: \(\mathcal{B} \subseteq \mathcal{S}\).

By 2.29 it suffices to show that every open subset of \(\mathbb{R}\) belongs to \(\mathcal{S}\). Let \(G \subseteq \mathbb{R}\) be open and put

\begin{equation*} P = \{(p,q) \in \mathbb{Q} \times \mathbb{Q} : p < q \text{ and } (p,q) \subseteq G\}. \end{equation*}

We claim \(G = \bigcup_{(p,q) \in P} (p,q)\). The inclusion \(\supseteq\) holds by the definition of \(P\). For \(\subseteq\), let \(x \in G\); openness of \(G\) gives \(\varepsilon > 0\) with \((x-\varepsilon, x+\varepsilon) \subseteq G\), and by density of \(\mathbb{Q}\) we may choose rationals \(p \in (x - \varepsilon, x)\) and \(q \in (x, x + \varepsilon)\). Then \(x \in (p,q) \subseteq (x-\varepsilon, x+\varepsilon) \subseteq G\), so \((p,q) \in P\) and \(x\) lies in the right side. This proves the claim.

Since \(P \subseteq \mathbb{Q} \times \mathbb{Q}\) is countable, the displayed union is a countable union of sets that lie in \(\mathcal{S}\) by Step 2 (if \(P = \emptyset\) then \(G = \emptyset \in \mathcal{S}\)). Hence \(G \in \mathcal{S}\).

Thus \(\mathcal{S}\) contains every open subset of \(\mathbb{R}\), so \(\mathcal{B} \subseteq \mathcal{S}\). With Step 1 this gives \(\mathcal{S} = \mathcal{B}\), so \(\mathcal{S}\) is the collection of Borel subsets of \(\mathbb{R}\).

Problem (2B.7)

Prove that the collection of Borel subsets of \(\mathbb{R}\) is translation invariant. More precisely, prove that if \(B \subseteq \mathbb{R}\) is a Borel set and \(t \in \mathbb{R}\), then \(t + B\) is a Borel set.

Solution

Here \(t + B = \{t + b : b \in B\}\). Let \(\mathcal{B}\) denote the collection of Borel subsets of \(\mathbb{R}\), and fix \(t \in \mathbb{R}\).

A short proof. Define \(f : \mathbb{R} \to \mathbb{R}\) by \(f(x) = x - t\). Then \(f\) is continuous, so \(f\) is a Borel measurable function by 2.41. For \(B \subseteq \mathbb{R}\) we have

\begin{equation*} f^{-1}(B) = \{x \in \mathbb{R} : x - t \in B\} = \{t + b : b \in B\} = t + B, \end{equation*}

where the middle equality holds because \(x - t \in B\) if and only if \(x = t + b\) for some \(b \in B\). If \(B\) is a Borel set, then Borel measurability of \(f\) (see 2.40) gives that \(f^{-1}(B) = t + B\) is a Borel set.

A self-contained proof, not using 2.41. Define

\begin{equation*} \mathcal{T} = \{A \subseteq \mathbb{R} : t + A \in \mathcal{B}\}. \end{equation*}

We show that \(\mathcal{T}\) is a \(\sigma\)-algebra on \(\mathbb{R}\) that contains every open subset of \(\mathbb{R}\); it then follows from 2.29 that \(\mathcal{B} \subseteq \mathcal{T}\), which is exactly the assertion that \(t + B\) is a Borel set for every Borel set \(B\).

First, \(t + \emptyset = \emptyset \in \mathcal{B}\), so \(\emptyset \in \mathcal{T}\).

Second, for every \(A \subseteq \mathbb{R}\),

\begin{equation*} t + (\mathbb{R} \setminus A) = \mathbb{R} \setminus (t + A). \end{equation*}

Indeed, for \(x \in \mathbb{R}\) we have \(x \in t + A\) if and only if \(x - t \in A\); hence \(x \in t + (\mathbb{R} \setminus A)\) if and only if \(x - t \notin A\), if and only if \(x \notin t + A\). So if \(A \in \mathcal{T}\), then \(t + A \in \mathcal{B}\), and since \(\mathcal{B}\) is closed under complementation we get \(t + (\mathbb{R} \setminus A) = \mathbb{R} \setminus (t+A) \in \mathcal{B}\), that is, \(\mathbb{R} \setminus A \in \mathcal{T}\).

Third, if \(A_1, A_2, \ldots\) is a sequence of subsets of \(\mathbb{R}\), then

\begin{equation*} t + \bigcup_{k=1}^\infty A_k = \bigcup_{k=1}^\infty (t + A_k), \end{equation*}

because \(x - t \in \bigcup_k A_k\) if and only if \(x - t \in A_k\) for some \(k\). So if each \(A_k \in \mathcal{T}\), then each \(t + A_k \in \mathcal{B}\), and \(\mathcal{B}\) is closed under countable unions, giving \(\bigcup_{k=1}^\infty A_k \in \mathcal{T}\).

Thus \(\mathcal{T}\) is a \(\sigma\)-algebra on \(\mathbb{R}\).

Finally, suppose \(G \subseteq \mathbb{R}\) is open; we check \(t + G\) is open. Let \(x \in t + G\), so \(x - t \in G\). Since \(G\) is open there is \(\varepsilon > 0\) with \((x - t - \varepsilon,\ x - t + \varepsilon) \subseteq G\). If \(|y - x| < \varepsilon\) then \(|(y - t) - (x - t)| < \varepsilon\), so \(y - t \in G\) and hence \(y \in t + G\). Thus \((x - \varepsilon, x + \varepsilon) \subseteq t + G\), proving \(t + G\) is open. Every open set is a Borel set, so \(t + G \in \mathcal{B}\) and therefore \(G \in \mathcal{T}\).

So \(\mathcal{T}\) is a \(\sigma\)-algebra on \(\mathbb{R}\) containing all open subsets of \(\mathbb{R}\). Because \(\mathcal{B}\) is the smallest such \(\sigma\)-algebra (2.29), we conclude \(\mathcal{B} \subseteq \mathcal{T}\). In other words, if \(B\) is a Borel subset of \(\mathbb{R}\) and \(t \in \mathbb{R}\), then \(t + B\) is a Borel set. Hence the collection of Borel subsets of \(\mathbb{R}\) is translation invariant.

Problem (2B.8)

Prove that the collection of Borel subsets of \(\mathbb{R}\) is dilation invariant. More precisely, prove that if \(B \subseteq \mathbb{R}\) is a Borel set and \(t \in \mathbb{R}\), then \(tB\) (which is defined to be \(\{tb : b \in B\}\)) is a Borel set.

Solution

Suppose \(B \subseteq \mathbb{R}\) is a Borel set and \(t \in \mathbb{R}\).

First dispose of the degenerate case \(t = 0\). If \(B = \emptyset\) then \(0B = \emptyset\), and if \(B \neq \emptyset\) then \(0B = \{0\}\). Both \(\emptyset\) and \(\{0\}\) are Borel sets (\(\{0\}\) is closed, hence Borel by 2.30), so \(tB\) is a Borel set when \(t = 0\).

Now suppose \(t \neq 0\). Define \(m : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} m(x) = \frac{x}{t}. \end{equation*}

Then \(m\) is continuous, and \(\mathbb{R}\) is a Borel set, so 2.41 tells us that \(m\) is a Borel measurable function. Hence \(m^{-1}(B)\) is a Borel set, by the definition 2.40 of Borel measurability.

We now check that \(m^{-1}(B) = tB\). For \(x \in \mathbb{R}\) we have, using 2.31,

\begin{equation*} x \in m^{-1}(B) \iff \frac{x}{t} \in B . \end{equation*}

If \(x/t \in B\), then \(x = t\,(x/t) \in tB\). Conversely, if \(x \in tB\), then \(x = tb\) for some \(b \in B\), and since \(t \neq 0\) this gives \(x/t = b \in B\). Thus

\begin{equation*} m^{-1}(B) = tB, \end{equation*}

and therefore \(tB\) is a Borel set.

Combining the two cases: \(tB\) is a Borel set for every Borel set \(B \subseteq \mathbb{R}\) and every \(t \in \mathbb{R}\). In other words, the collection of Borel subsets of \(\mathbb{R}\) is dilation invariant.

It is worth noting the alternative argument, which avoids 2.41 and mirrors the proof of Exercise 7 in this section. Fix \(t \neq 0\) and let

\begin{equation*} \mathcal{S} = \{ B \subseteq \mathbb{R} : tB \text{ is a Borel set} \}. \end{equation*}

Because \(t\emptyset = \emptyset\), we have \(\emptyset \in \mathcal{S}\). Because multiplication by \(t\) is a bijection of \(\mathbb{R}\) onto \(\mathbb{R}\), we have \(t(\mathbb{R} \setminus B) = \mathbb{R} \setminus tB\) and \(t\bigl(\bigcup_{k=1}^{\infty} B_k\bigr) = \bigcup_{k=1}^{\infty} tB_k\); hence \(\mathcal{S}\) is closed under complementation and under countable unions, so \(\mathcal{S}\) is a \(\sigma\)-algebra on \(\mathbb{R}\). If \(G \subseteq \mathbb{R}\) is open, then \(tG\) is open (multiplication by \(t\) is a homeomorphism) and thus Borel, so \(G \in \mathcal{S}\). Hence \(\mathcal{S}\) is a \(\sigma\)-algebra containing every open subset of \(\mathbb{R}\), and so \(\mathcal{S}\) contains the smallest such \(\sigma\)-algebra, which by 2.29 is the collection of Borel sets. Thus \(tB\) is a Borel set for every Borel set \(B\).

Problem (2B.9)

Give an example of a measurable space \((X, \mathcal{S})\) and a function \(f : X \to \mathbb{R}\) such that \(|f|\) is \(\mathcal{S}\)-measurable but \(f\) is not \(\mathcal{S}\)-measurable.

Solution

Take \(X = \mathbb{R}\) and \(\mathcal{S} = \{\emptyset, \mathbb{R}\}\). This is a \(\sigma\)-algebra on \(\mathbb{R}\) (it contains \(\emptyset\), it is closed under complementation because \(\mathbb{R} \setminus \emptyset = \mathbb{R}\) and \(\mathbb{R} \setminus \mathbb{R} = \emptyset\), and any union of sets from \(\mathcal{S}\) is either \(\emptyset\) or \(\mathbb{R}\)), so \((\mathbb{R}, \mathcal{S})\) is a measurable space.

Define \(f : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} f(x) = \begin{cases} 1 & \text{if } x \geq 0,\\ -1 & \text{if } x < 0. \end{cases} \end{equation*}

Then \(|f|\) is the constant function with value \(1\). For every Borel set \(B \subseteq \mathbb{R}\) we have \(|f|^{-1}(B) = \mathbb{R}\) if \(1 \in B\) and \(|f|^{-1}(B) = \emptyset\) if \(1 \notin B\); in either case \(|f|^{-1}(B) \in \mathcal{S}\). Hence \(|f|\) is \(\mathcal{S}\)-measurable. (This is also the content of the first bullet point of 2.36: when \(\mathcal{S} = \{\emptyset, X\}\), the \(\mathcal{S}\)-measurable functions are exactly the constant functions.)

However \(f\) is not \(\mathcal{S}\)-measurable. Indeed, \((0, \infty)\) is an open, hence Borel, subset of \(\mathbb{R}\), and

\begin{equation*} f^{-1}\bigl((0,\infty)\bigr) = \{x \in \mathbb{R} : f(x) > 0\} = [0, \infty), \end{equation*}

which is neither \(\emptyset\) nor \(\mathbb{R}\), so \(f^{-1}\bigl((0,\infty)\bigr) \notin \mathcal{S}\). Thus \(f\) fails the definition 2.35 of \(\mathcal{S}\)-measurability.

The example shows that the implication in 2.45 (if \(f\) is measurable then \(|f|\) is measurable) cannot be reversed: passing to \(|f|\) destroys exactly the sign information that measurability of \(f\) would require.

More generally, the same phenomenon occurs for any measurable space \((X, \mathcal{S})\) and any set \(E \subseteq X\) with \(E \notin \mathcal{S}\): the function \(f = \chi_E - \chi_{X \setminus E}\) satisfies \(|f| \equiv 1\), which is \(\mathcal{S}\)-measurable, while \(f^{-1}\bigl((0,\infty)\bigr) = E \notin \mathcal{S}\), so \(f\) is not \(\mathcal{S}\)-measurable. Such an \(E\) exists whenever \(\mathcal{S}\) is not the collection of all subsets of \(X\).

Problem (2B.10)

Show that the set of real numbers that have a decimal expansion with the digit 5 appearing infinitely often is a Borel set.

Solution

Let \(A\) denote the set in question. Here a decimal expansion of \(x \geq 0\) means a representation

\begin{equation*} x = N + \sum_{j=1}^{\infty} e_j 10^{-j}, \qquad N \in \mathbb{Z},\ N \ge 0,\ e_j \in \{0, 1, \ldots, 9\}, \end{equation*}

written \(x = N.e_1 e_2 e_3 \ldots\), and a decimal expansion of a negative number is a minus sign followed by a decimal expansion of its absolute value. Since the digits of the integer part \(N\) form a finite string, the digit 5 can appear infinitely often only among the digits \(e_1, e_2, \ldots\) after the decimal point.

Step 1: reduction to \([0, \infty)\). Because the expansions of \(x\) and of \(-x\) carry the same digit string, \(x \in A\) if and only if \(|x| \in A\). Hence, setting \(A_+ = A \cap [0, \infty)\), we have

\begin{equation*} A = A_+ \cup (-A_+), \end{equation*}

where \(-A_+ = \{-y : y \in A_+\}\). By Exercise 8 of this section (dilation invariance, with \(t = -1\)), \(-A_+\) is a Borel set whenever \(A_+\) is. So it suffices to prove that \(A_+\) is a Borel set.

Step 2: the canonical digits. For \(x \geq 0\) and \(k \in \mathbb{Z}^+\) define

\begin{equation*} d_k(x) = \lfloor 10^k x \rfloor - 10 \lfloor 10^{k-1} x \rfloor . \end{equation*}

Write \(q = \lfloor 10^{k-1} x \rfloor\), so \(q \leq 10^{k-1} x < q + 1\) and hence \(10q \leq 10^k x < 10q + 10\), which gives \(10q \leq \lfloor 10^k x \rfloor \leq 10q + 9\). Thus \(d_k(x) \in \{0, 1, \ldots, 9\}\). Moreover the sum telescopes:

\begin{equation*} \lfloor x \rfloor + \sum_{k=1}^{n} d_k(x) 10^{-k} = \frac{\lfloor 10^n x \rfloor}{10^n}, \end{equation*}

and \(\lfloor 10^n x\rfloor / 10^n \to x\) as \(n \to \infty\) because \(x - 10^{-n} < \lfloor 10^n x \rfloor/10^n \le x\). Hence \(\lfloor x \rfloor . d_1(x) d_2(x) d_3(x)\ldots\) is a decimal expansion of \(x\); call it the canonical expansion.

Step 3: 5 appears infinitely often in some expansion if and only if it appears infinitely often in the canonical one. One direction is immediate, since the canonical string is itself an expansion. For the other direction, suppose \(x = N.e_1e_2e_3\ldots\) is a decimal expansion of \(x\) in which \(e_j = 5\) for infinitely many \(j\). Fix \(k \geq 0\). Since \(e_j = 5\) for infinitely many \(j\), there is some \(i > k\) with \(e_i = 5\); in particular the digits \(e_{k+1}, e_{k+2}, \ldots\) are not all equal to \(9\). Hence

\begin{equation*} \sum_{j > k} e_j 10^{-j} \leq \Bigl( \sum_{j > k} 9 \cdot 10^{-j} \Bigr) - 4 \cdot 10^{-i} = 10^{-k} - 4 \cdot 10^{-i} < 10^{-k}, \end{equation*}

and of course \(\sum_{j>k} e_j 10^{-j} \geq 0\). Setting \(M_k = 10^k N + \sum_{j=1}^{k} e_j 10^{k-j} \in \mathbb{Z}\), the displayed inequality gives

\begin{equation*} M_k \leq 10^k x < M_k + 1, \end{equation*}

so \(\lfloor 10^k x \rfloor = M_k\) for every \(k \geq 0\). Consequently \(\lfloor x \rfloor = M_0 = N\) and, for \(k \in \mathbb{Z}^+\),

\begin{equation*} d_k(x) = M_k - 10 M_{k-1} = e_k . \end{equation*}

So the given expansion is the canonical one, and therefore \(d_k(x) = 5\) for infinitely many \(k\). Hence

\begin{equation*} A_+ = \{ x \geq 0 : d_k(x) = 5 \text{ for infinitely many } k \in \mathbb{Z}^+ \}. \end{equation*}

Step 4: each digit set is Borel. Fix \(k \in \mathbb{Z}^+\) and let

\begin{equation*} D_k = \{ x \geq 0 : d_k(x) = 5 \}. \end{equation*}

We claim \(d_k(x) = 5\) if and only if \(\lfloor 10^k x \rfloor \in \{10m + 5 : m \in \mathbb{Z}\}\). Indeed, if \(d_k(x) = 5\) then \(\lfloor 10^k x\rfloor = 10\lfloor 10^{k-1}x\rfloor + 5\), which has the stated form. Conversely, if \(\lfloor 10^k x \rfloor = 10m + 5\) for some \(m \in \mathbb{Z}\), then \(10m + 5 \leq 10^k x < 10m + 6\), so \(m + \tfrac12 \leq 10^{k-1} x < m + \tfrac35\), whence \(\lfloor 10^{k-1} x \rfloor = m\) and \(d_k(x) = (10m + 5) - 10m = 5\). Since \(x \geq 0\) forces \(\lfloor 10^k x\rfloor \ge 0\), only \(m \geq 0\) occurs, and the claim translates into

\begin{equation*} D_k = \bigcup_{m=0}^{\infty} \left[ \frac{10m + 5}{10^{k}},\ \frac{10m + 6}{10^{k}} \right). \end{equation*}

Each half-open interval \([a, b)\) is a Borel set by 2.30, and \(D_k\) is a countable union of such intervals, so \(D_k\) is a Borel set.

Step 5: assembling. The digit 5 occurs among \(d_1(x), d_2(x), \ldots\) infinitely often exactly when for every \(n \in \mathbb{Z}^+\) there is some \(k \geq n\) with \(d_k(x) = 5\). Therefore

\begin{equation*} A_+ = \bigcap_{n=1}^{\infty} \bigcup_{k=n}^{\infty} D_k . \end{equation*}

Each inner union is a countable union of Borel sets, hence Borel by the definition 2.23 of a \(\sigma\)-algebra, and the outer intersection is a countable intersection of Borel sets, hence Borel by 2.25(c). Thus \(A_+\) is a Borel set, and by Step 1 so is

\begin{equation*} A = A_+ \cup (-A_+). \end{equation*}

Problem (2B.11)

Suppose \(\mathcal{T}\) is a \(\sigma\)-algebra on a set \(Y\) and \(X \in \mathcal{T}\). Let \(\mathcal{S} = \{E \in \mathcal{T} : E \subseteq X\}\).

(a) Show that \(\mathcal{S} = \{F \cap X : F \in \mathcal{T}\}\).

(b) Show that \(\mathcal{S}\) is a \(\sigma\)-algebra on \(X\).

Solution

(a) First suppose \(E \in \mathcal{S}\). Then \(E \in \mathcal{T}\) and \(E \subseteq X\), so \(E = E \cap X\), which exhibits \(E\) as an element of \(\{F \cap X : F \in \mathcal{T}\}\) (take \(F = E\)). Thus

\begin{equation*} \mathcal{S} \subseteq \{F \cap X : F \in \mathcal{T}\}. \end{equation*}

For the other inclusion, suppose \(F \in \mathcal{T}\). Since \(X \in \mathcal{T}\) by hypothesis, 2.25(b) gives \(F \cap X \in \mathcal{T}\). Also \(F \cap X \subseteq X\). Hence \(F \cap X \in \mathcal{S}\), proving

\begin{equation*} \{F \cap X : F \in \mathcal{T}\} \subseteq \mathcal{S}. \end{equation*}

The two inclusions together give \(\mathcal{S} = \{F \cap X : F \in \mathcal{T}\}\).

(b) We verify the three conditions in the definition 2.23 of a \(\sigma\)-algebra on \(X\). Note first that every element of \(\mathcal{S}\) is by definition a subset of \(X\), so \(\mathcal{S}\) is indeed a collection of subsets of \(X\); and complementation below must be taken within \(X\), not within \(Y\).

The empty set: \(\emptyset \in \mathcal{T}\) because \(\mathcal{T}\) is a \(\sigma\)-algebra, and \(\emptyset \subseteq X\). Hence \(\emptyset \in \mathcal{S}\).

Closure under complementation in \(X\): suppose \(E \in \mathcal{S}\). Then \(E \in \mathcal{T}\) and \(X \in \mathcal{T}\), so \(X \setminus E \in \mathcal{T}\) by 2.25(b). Since also \(X \setminus E \subseteq X\), we conclude \(X \setminus E \in \mathcal{S}\).

Closure under countable unions: suppose \(E_1, E_2, \ldots\) is a sequence of elements of \(\mathcal{S}\). Each \(E_k \in \mathcal{T}\), so \(\bigcup_{k=1}^{\infty} E_k \in \mathcal{T}\) because \(\mathcal{T}\) is a \(\sigma\)-algebra. Each \(E_k \subseteq X\), so \(\bigcup_{k=1}^{\infty} E_k \subseteq X\). Hence \(\bigcup_{k=1}^{\infty} E_k \in \mathcal{S}\).

Thus \(\mathcal{S}\) is a \(\sigma\)-algebra on \(X\).

The point of this exercise is recorded in the text just before 2.40: if \(X\) is a Borel subset of \(\mathbb{R}\) and \(\mathcal{T}\) is the collection of Borel subsets of \(\mathbb{R}\), then \(\mathcal{S}\) is the collection of Borel sets contained in \(X\), and part (a) says that this is the same as the collection of sets of the form \(F \cap X\) with \(F\) Borel. This is the \(\sigma\)-algebra with respect to which Borel measurable functions on \(X\) (definition 2.40) are measurable.

Problem (2B.12)

Suppose \(f : \mathbb{R} \to \mathbb{R}\) is a function.

(a) For \(k \in \mathbb{Z}^+\), let

\begin{equation*} G_k = \Bigl\{ a \in \mathbb{R} : \text{there exists } \delta > 0 \text{ such that } |f(b) - f( c)| < \tfrac{1}{k} \text{ for all } b, c \in (a - \delta, a + \delta) \Bigr\}. \end{equation*}

Prove that \(G_k\) is an open subset of \(\mathbb{R}\) for each \(k \in \mathbb{Z}^+\).

(b) Prove that the set of points at which \(f\) is continuous equals \(\bigcap_{k=1}^{\infty} G_k\).

(c) Conclude that the set of points at which \(f\) is continuous is a Borel set.

Solution

(a) Fix \(k \in \mathbb{Z}^+\) and suppose \(a \in G_k\). Choose \(\delta > 0\) such that

\begin{equation*} |f(b) - f( c)| < \tfrac{1}{k} \quad \text{for all } b, c \in (a - \delta, a + \delta). \end{equation*}

We claim that \((a - \tfrac{\delta}{2}, a + \tfrac{\delta}{2}) \subseteq G_k\), which shows that \(G_k\) contains an open interval about each of its points and hence is open.

To prove the claim, suppose \(a^{\prime} \in (a - \tfrac{\delta}{2}, a + \tfrac{\delta}{2})\) and set

\begin{equation*} \delta^{\prime} = \tfrac{\delta}{2} - |a^{\prime} - a| > 0 . \end{equation*}

If \(y \in (a^{\prime} - \delta^{\prime}, a^{\prime} + \delta^{\prime})\), then

\begin{equation*} |y - a| \leq |y - a^{\prime}| + |a^{\prime} - a| < \delta^{\prime} + |a^{\prime} - a| = \tfrac{\delta}{2} < \delta, \end{equation*}

so \(y \in (a - \delta, a + \delta)\). Thus \((a^{\prime} - \delta^{\prime}, a^{\prime} + \delta^{\prime}) \subseteq (a - \delta, a + \delta)\), and therefore \(|f(b) - f( c)| < \tfrac{1}{k}\) for all \(b, c \in (a^{\prime} - \delta^{\prime}, a^{\prime} + \delta^{\prime})\). Hence \(a^{\prime} \in G_k\), proving the claim and (a).

(b) Let \(C\) denote the set of points at which \(f\) is continuous.

Suppose \(a \in C\) and let \(k \in \mathbb{Z}^+\). By continuity of \(f\) at \(a\) applied with \(\varepsilon = \tfrac{1}{2k}\), there exists \(\delta > 0\) such that \(|f(y) - f(a)| < \tfrac{1}{2k}\) for all \(y \in (a - \delta, a + \delta)\). Hence for all \(b, c \in (a - \delta, a + \delta)\),

\begin{equation*} |f(b) - f( c)| \leq |f(b) - f(a)| + |f(a) - f( c)| < \tfrac{1}{2k} + \tfrac{1}{2k} = \tfrac{1}{k}. \end{equation*}

Thus \(a \in G_k\). Since \(k\) was arbitrary, \(a \in \bigcap_{k=1}^{\infty} G_k\), and so \(C \subseteq \bigcap_{k=1}^{\infty} G_k\).

Conversely, suppose \(a \in \bigcap_{k=1}^{\infty} G_k\) and let \(\varepsilon > 0\). Choose \(k \in \mathbb{Z}^+\) with \(\tfrac{1}{k} < \varepsilon\). Because \(a \in G_k\), there exists \(\delta > 0\) such that \(|f(b) - f( c)| < \tfrac{1}{k}\) for all \(b, c \in (a - \delta, a + \delta)\). Taking \(c = a\) (which lies in \((a - \delta, a + \delta)\)), we get

\begin{equation*} |f(b) - f(a)| < \tfrac{1}{k} < \varepsilon \quad \text{for all } b \in (a - \delta, a + \delta). \end{equation*}

Hence \(f\) is continuous at \(a\), so \(a \in C\). This gives \(\bigcap_{k=1}^{\infty} G_k \subseteq C\).

The two inclusions show that \(C = \bigcap_{k=1}^{\infty} G_k\).

(c) By (a), each \(G_k\) is an open subset of \(\mathbb{R}\) and hence is a Borel set, because by 2.29 the collection of Borel sets is a \(\sigma\)-algebra containing all open subsets of \(\mathbb{R}\). By 2.25(c) a \(\sigma\)-algebra is closed under countable intersections, so

\begin{equation*} C = \bigcap_{k=1}^{\infty} G_k \end{equation*}

is a Borel set. Thus the set of points at which \(f\) is continuous is a Borel set, as asserted in the last bullet point of 2.30.

Problem (2B.13)

Suppose \((X, \mathcal{S})\) is a measurable space, \(E_1, \ldots, E_n\) are disjoint subsets of \(X\), and \(c_1, \ldots, c_n\) are distinct nonzero real numbers. Prove that \(c_1 \chi_{E_1} + \cdots + c_n \chi_{E_n}\) is an \(\mathcal{S}\)-measurable function if and only if \(E_1, \ldots, E_n \in \mathcal{S}\).

Solution

Write

\begin{equation*} f = c_1 \chi_{E_1} + \cdots + c_n \chi_{E_n} . \end{equation*}

First suppose \(E_1, \ldots, E_n \in \mathcal{S}\). Fix \(j\). By 2.38, for every \(B \subseteq \mathbb{R}\) the inverse image \(\chi_{E_j}^{-1}(B)\) equals one of \(E_j\), \(X \setminus E_j\), \(X\), \(\emptyset\), according to which of \(0, 1\) belong to \(B\). All four of these sets are in \(\mathcal{S}\): we have \(E_j \in \mathcal{S}\) by hypothesis, \(X \in \mathcal{S}\) and \(\emptyset \in \mathcal{S}\) by 2.25(a) and the definition 2.23, and \(X \setminus E_j \in \mathcal{S}\) by the second bullet point of 2.23. Hence \(\chi_{E_j}^{-1}(B) \in \mathcal{S}\) for every Borel set \(B \subseteq \mathbb{R}\), so \(\chi_{E_j}\) is an \(\mathcal{S}\)-measurable function. Each constant function on \(X\) is \(\mathcal{S}\)-measurable (the inverse image of a set is either \(X\) or \(\emptyset\), both of which are in \(\mathcal{S}\) by 2.25(a)), so 2.46(a) implies that each product \(c_j \chi_{E_j}\) is \(\mathcal{S}\)-measurable. Applying 2.46(a) again, this time to sums, and inducting on the number of summands, we conclude that \(f\) is an \(\mathcal{S}\)-measurable function.

Now suppose \(f\) is \(\mathcal{S}\)-measurable. We claim that

\begin{equation*} f^{-1}(\{c_j\}) = E_j \quad \text{for each } j \in \{1, \ldots, n\}. \end{equation*}

To see this, first note the values of \(f\). If \(x \in E_i\) for some \(i\), then \(\chi_{E_i}(x) = 1\) and, because \(E_1, \ldots, E_n\) are disjoint, \(\chi_{E_m}(x) = 0\) for every \(m \neq i\); hence \(f(x) = c_i\). If \(x \notin E_1 \cup \cdots \cup E_n\), then every term vanishes and \(f(x) = 0\).

Fix \(j\). If \(x \in E_j\), then \(f(x) = c_j\) by the previous paragraph, so \(x \in f^{-1}(\{c_j\})\). Conversely suppose \(f(x) = c_j\). If \(x \notin E_1 \cup \cdots \cup E_n\), then \(f(x) = 0 \neq c_j\), because the numbers \(c_1, \ldots, c_n\) are nonzero; this is a contradiction. Hence \(x \in E_i\) for some \(i\), and then \(c_i = f(x) = c_j\), which forces \(i = j\) because the numbers \(c_1, \ldots, c_n\) are distinct. Thus \(x \in E_j\). This proves the claim.

The set \(\{c_j\}\) is a closed subset of \(\mathbb{R}\) and hence is a Borel set by 2.30. Because \(f\) is \(\mathcal{S}\)-measurable, the definition 2.35 gives

\begin{equation*} E_j = f^{-1}(\{c_j\}) \in \mathcal{S}. \end{equation*}

Since \(j \in \{1, \ldots, n\}\) was arbitrary, \(E_1, \ldots, E_n \in \mathcal{S}\).

Both hypotheses on the \(c_j\) are needed for the forward implication. If some \(c_j\) were \(0\), the corresponding term would contribute nothing and \(E_j\) could be an arbitrary (possibly non-measurable) set. If two of the \(c_j\) coincided, say \(c_1 = c_2\), then \(f = c_1 \chi_{E_1 \cup E_2} + \cdots\) would only detect \(E_1 \cup E_2\), which can be measurable while \(E_1\) and \(E_2\) are not. Disjointness is what makes \(f\) take the value \(c_j\) exactly on \(E_j\).

Problem (2B.14)

(a) Suppose \(f_1, f_2, \ldots\) is a sequence of functions from a set \(X\) to \(\mathbb{R}\). Explain why

\begin{equation*} \{x \in X : \text{the sequence } f_1(x), f_2(x), \ldots \text{ has a limit in } \mathbb{R}\} = \bigcap_{n=1}^{\infty} \bigcup_{j=1}^{\infty} \bigcap_{k=j}^{\infty} (f_j - f_k)^{-1}\Bigl( \bigl(-\tfrac{1}{n}, \tfrac{1}{n}\bigr) \Bigr). \end{equation*}

(b) Suppose \((X, \mathcal{S})\) is a measurable space and \(f_1, f_2, \ldots\) is a sequence of \(\mathcal{S}\)-measurable functions from \(X\) to \(\mathbb{R}\). Prove that

\begin{equation*} \{x \in X : \text{the sequence } f_1(x), f_2(x), \ldots \text{ has a limit in } \mathbb{R}\} \end{equation*}

is an \(\mathcal{S}\)-measurable subset of \(X\).

Solution

(a) Denote the left side by \(L\) and the right side by \(R\).

We first unwind the right side. By the definition 2.31 of inverse image, for \(x \in X\) we have

\begin{equation*} x \in (f_j - f_k)^{-1}\Bigl( \bigl(-\tfrac{1}{n}, \tfrac{1}{n}\bigr) \Bigr) \iff -\tfrac{1}{n} < f_j(x) - f_k(x) < \tfrac{1}{n} \iff |f_j(x) - f_k(x)| < \tfrac{1}{n}. \end{equation*}

Hence, reading the intersections as “for all” and the union as “there exists”,

\begin{equation*} x \in R \iff \text{for every } n \in \mathbb{Z}^+ \text{ there exists } j \in \mathbb{Z}^+ \text{ such that } |f_j(x) - f_k(x)| < \tfrac{1}{n} \text{ for all } k \geq j. \end{equation*}

Now fix \(x \in X\) and consider the sequence of real numbers \(f_1(x), f_2(x), \ldots\).

Suppose \(x \in L\), so the sequence has a limit in \(\mathbb{R}\); then it is a Cauchy sequence. Let \(n \in \mathbb{Z}^+\). By the Cauchy condition applied with \(\varepsilon = \tfrac{1}{n}\), there exists \(j \in \mathbb{Z}^+\) such that \(|f_p(x) - f_q(x)| < \tfrac{1}{n}\) for all \(p, q \geq j\). In particular, taking \(p = j\), we get \(|f_j(x) - f_k(x)| < \tfrac{1}{n}\) for all \(k \geq j\). Since \(n\) was arbitrary, \(x \in R\). Thus \(L \subseteq R\).

Conversely, suppose \(x \in R\), and let \(\varepsilon > 0\). Choose \(n \in \mathbb{Z}^+\) with \(\tfrac{2}{n} < \varepsilon\). By the description of \(R\) above, there exists \(j \in \mathbb{Z}^+\) such that \(|f_j(x) - f_k(x)| < \tfrac{1}{n}\) for all \(k \geq j\). Hence for all \(p, q \geq j\),

\begin{equation*} |f_p(x) - f_q(x)| \leq |f_p(x) - f_j(x)| + |f_j(x) - f_q(x)| < \tfrac{1}{n} + \tfrac{1}{n} = \tfrac{2}{n} < \varepsilon . \end{equation*}

Thus \(f_1(x), f_2(x), \ldots\) is a Cauchy sequence of real numbers, and so by the completeness of \(\mathbb{R}\) it has a limit in \(\mathbb{R}\). Hence \(x \in L\), giving \(R \subseteq L\).

The two inclusions establish the asserted equality. In words: the identity is exactly the statement that convergence in \(\mathbb{R}\) is equivalent to the Cauchy condition, with the Cauchy condition written using only the single “anchor” index \(j\) rather than two free indices.

(b) Let

\begin{equation*} L = \{x \in X : \text{the sequence } f_1(x), f_2(x), \ldots \text{ has a limit in } \mathbb{R}\}. \end{equation*}

Fix \(j, k \in \mathbb{Z}^+\). Since \(f_j\) and \(f_k\) are \(\mathcal{S}\)-measurable, 2.46(a) implies that \(f_j - f_k\) is an \(\mathcal{S}\)-measurable function from \(X\) to \(\mathbb{R}\). For each \(n \in \mathbb{Z}^+\) the interval \(\bigl(-\tfrac1n, \tfrac1n\bigr)\) is an open subset of \(\mathbb{R}\) and hence a Borel set, by 2.29. Therefore the definition 2.35 of measurability gives

\begin{equation*} (f_j - f_k)^{-1}\Bigl( \bigl(-\tfrac{1}{n}, \tfrac{1}{n}\bigr) \Bigr) \in \mathcal{S} \end{equation*}

for all \(n, j, k \in \mathbb{Z}^+\).

Now build up the set in (a) from these pieces. For fixed \(n\) and \(j\), the set

\begin{equation*} \bigcap_{k=j}^{\infty} (f_j - f_k)^{-1}\Bigl( \bigl(-\tfrac{1}{n}, \tfrac{1}{n}\bigr) \Bigr) \end{equation*}

is a countable intersection of elements of \(\mathcal{S}\), hence is in \(\mathcal{S}\) by 2.25(c). Next, for fixed \(n\), the set

\begin{equation*} \bigcup_{j=1}^{\infty} \bigcap_{k=j}^{\infty} (f_j - f_k)^{-1}\Bigl( \bigl(-\tfrac{1}{n}, \tfrac{1}{n}\bigr) \Bigr) \end{equation*}

is a countable union of elements of \(\mathcal{S}\), hence is in \(\mathcal{S}\) by the third bullet point in the definition 2.23 of a \(\sigma\)-algebra. Finally,

\begin{equation*} L = \bigcap_{n=1}^{\infty} \Biggl( \bigcup_{j=1}^{\infty} \bigcap_{k=j}^{\infty} (f_j - f_k)^{-1}\Bigl( \bigl(-\tfrac{1}{n}, \tfrac{1}{n}\bigr) \Bigr) \Biggr) \end{equation*}

is a countable intersection of elements of \(\mathcal{S}\), hence \(L \in \mathcal{S}\) by 2.25(c) again.

Thus the set of \(x \in X\) at which the sequence \(f_1(x), f_2(x), \ldots\) converges in \(\mathbb{R}\) is an \(\mathcal{S}\)-measurable subset of \(X\).

This complements 2.48: there the sequence was assumed to converge at every point of \(X\), and the limit function was shown to be \(\mathcal{S}\)-measurable. Here we see that even without that assumption the convergence set \(L\) is itself measurable. Consequently one may define \(f : X \to \mathbb{R}\) by \(f(x) = \lim_{k \to \infty} f_k(x)\) for \(x \in L\) and \(f(x) = 0\) for \(x \in X \setminus L\), and this \(f\) is \(\mathcal{S}\)-measurable. Indeed, let \(\mathcal{S}_L = \{E \in \mathcal{S} : E \subseteq L\}\), which by Exercise 11 of this section is a \(\sigma\)-algebra on \(L\). Each restriction \(f_k|_L\) is \(\mathcal{S}_L\)-measurable, because \((f_k|_L)^{-1}(B) = f_k^{-1}(B) \cap L \in \mathcal{S}_L\) for every Borel set \(B\). Since \(f_1(x), f_2(x), \ldots\) converges for every \(x \in L\), 2.48 shows that \(f|_L\) is \(\mathcal{S}_L\)-measurable. Finally, for every Borel set \(B \subseteq \mathbb{R}\),

\begin{equation*} f^{-1}(B) = \begin{cases} (f|_L)^{-1}(B) \cup (X \setminus L) & \text{if } 0 \in B,\\ (f|_L)^{-1}(B) & \text{if } 0 \notin B, \end{cases} \end{equation*}

and both \((f|_L)^{-1}(B) \in \mathcal{S}_L \subseteq \mathcal{S}\) and \(X \setminus L \in \mathcal{S}\), so \(f^{-1}(B) \in \mathcal{S}\).

Problem (2B.15)

Suppose \(X\) is a set and \(E_1, E_2, \ldots\) is a disjoint sequence of subsets of \(X\) such that \(\bigcup_{k=1}^{\infty} E_k = X\). Let \(\mathcal{S} = \left\{ \bigcup_{k \in K} E_k : K \subseteq \mathbb{Z}^+ \right\}\).

(a) Show that \(\mathcal{S}\) is a \(\sigma\)-algebra on \(X\).

(b) Prove that a function from \(X\) to \(\mathbb{R}\) is \(\mathcal{S}\)-measurable if and only if the function is constant on \(E_k\) for every \(k \in \mathbb{Z}^+\).

Solution

Throughout, we use the following basic consequence of the hypotheses: because \(\bigcup_{k=1}^{\infty} E_k = X\) and the sets \(E_1, E_2, \ldots\) are pairwise disjoint, every \(x \in X\) belongs to \(E_k\) for exactly one \(k \in \mathbb{Z}^+\). Denote that index by \(k(x)\).

Part (a). Taking \(K = \emptyset\) gives \(\bigcup_{k \in \emptyset} E_k = \emptyset\), so \(\emptyset \in \mathcal{S}\).

Closure under complementation. Suppose \(E \in \mathcal{S}\), say \(E = \bigcup_{k \in K} E_k\) with \(K \subseteq \mathbb{Z}^+\). We claim

\begin{equation*} X \setminus E = \bigcup_{k \in \mathbb{Z}^+ \setminus K} E_k . \end{equation*}

Indeed, let \(x \in X\). Then \(x \in E\) if and only if \(x \in E_j\) for some \(j \in K\), which (by uniqueness of \(k(x)\)) happens if and only if \(k(x) \in K\). Hence \(x \in X \setminus E\) if and only if \(k(x) \in \mathbb{Z}^+ \setminus K\), which happens if and only if \(x \in \bigcup_{k \in \mathbb{Z}^+ \setminus K} E_k\). This proves the displayed equation, and therefore \(X \setminus E \in \mathcal{S}\).

Closure under countable unions. Suppose \(A_1, A_2, \ldots \in \mathcal{S}\), say \(A_n = \bigcup_{k \in K_n} E_k\) with \(K_n \subseteq \mathbb{Z}^+\). Put \(K = \bigcup_{n=1}^{\infty} K_n \subseteq \mathbb{Z}^+\). Then

\begin{equation*} \bigcup_{n=1}^{\infty} A_n = \bigcup_{n=1}^{\infty} \bigcup_{k \in K_n} E_k = \bigcup_{k \in K} E_k \in \mathcal{S} . \end{equation*}

Thus \(\mathcal{S}\) contains \(\emptyset\) and is closed under complementation and countable unions, so \(\mathcal{S}\) is a \(\sigma\)-algebra on \(X\).

Part (b). Suppose first that \(f : X \to \mathbb{R}\) is constant on \(E_k\) for every \(k \in \mathbb{Z}^+\). Let \(B \subseteq \mathbb{R}\) be a Borel set and set

\begin{equation*} K = \{ k \in \mathbb{Z}^+ : E_k \subseteq f^{-1}(B) \} . \end{equation*}

Clearly \(\bigcup_{k \in K} E_k \subseteq f^{-1}(B)\). Conversely, suppose \(x \in f^{-1}(B)\). Because \(f\) is constant on \(E_{k(x)}\) and \(x \in E_{k(x)}\), we have \(f(y) = f(x) \in B\) for every \(y \in E_{k(x)}\); hence \(E_{k(x)} \subseteq f^{-1}(B)\), so \(k(x) \in K\) and thus \(x \in \bigcup_{k \in K} E_k\). Therefore

\begin{equation*} f^{-1}(B) = \bigcup_{k \in K} E_k \in \mathcal{S} . \end{equation*}

Because \(B\) was an arbitrary Borel subset of \(\mathbb{R}\), the function \(f\) is \(\mathcal{S}\)-measurable.

Conversely, suppose \(f : X \to \mathbb{R}\) is \(\mathcal{S}\)-measurable. Fix \(k \in \mathbb{Z}^+\) and suppose \(x, y \in E_k\); we must show \(f(x) = f(y)\). The singleton \(\{f(x)\}\) is a closed subset of \(\mathbb{R}\) and hence a Borel set, so \(f^{-1}(\{f(x)\}) \in \mathcal{S}\); write \(f^{-1}(\{f(x)\}) = \bigcup_{j \in K} E_j\) for some \(K \subseteq \mathbb{Z}^+\). Now \(x \in f^{-1}(\{f(x)\})\), so \(x \in E_j\) for some \(j \in K\); by uniqueness of the index containing \(x\) we get \(j = k(x) = k\), so \(k \in K\). Consequently \(E_k \subseteq f^{-1}(\{f(x)\})\), and since \(y \in E_k\) this gives \(f(y) = f(x)\). Thus \(f\) is constant on \(E_k\) for every \(k \in \mathbb{Z}^+\).

Problem (2B.16)

Suppose \(\mathcal{S}\) is a \(\sigma\)-algebra on a set \(X\) and \(A \subseteq X\). Let

\begin{equation*} \mathcal{S}_A = \{ E \in \mathcal{S} : A \subseteq E \text{ or } A \cap E = \emptyset \}. \end{equation*}

(a) Prove that \(\mathcal{S}_A\) is a \(\sigma\)-algebra on \(X\).

(b) Suppose \(f : X \to \mathbb{R}\) is a function. Prove that \(f\) is measurable with respect to \(\mathcal{S}_A\) if and only if \(f\) is measurable with respect to \(\mathcal{S}\) and \(f\) is constant on \(A\).

Solution

Part (a). We have \(\emptyset \in \mathcal{S}\) and \(A \cap \emptyset = \emptyset\), so \(\emptyset \in \mathcal{S}_A\).

Closure under complementation. Suppose \(E \in \mathcal{S}_A\). Then \(E \in \mathcal{S}\), so \(X \setminus E \in \mathcal{S}\). If \(A \subseteq E\), then \(A \cap (X \setminus E) = \emptyset\). If instead \(A \cap E = \emptyset\), then every element of \(A\) lies in \(X\) but not in \(E\), so \(A \subseteq X \setminus E\). In either case \(X \setminus E \in \mathcal{S}_A\).

Closure under countable unions. Suppose \(E_1, E_2, \ldots \in \mathcal{S}_A\) and let \(E = \bigcup_{k=1}^{\infty} E_k\). Then \(E \in \mathcal{S}\) because \(\mathcal{S}\) is a \(\sigma\)-algebra. If \(A \subseteq E_k\) for some \(k\), then \(A \subseteq E\). Otherwise the definition of \(\mathcal{S}_A\) forces \(A \cap E_k = \emptyset\) for every \(k\), and then

\begin{equation*} A \cap E = A \cap \bigcup_{k=1}^{\infty} E_k = \bigcup_{k=1}^{\infty} (A \cap E_k) = \emptyset . \end{equation*}

In either case \(E \in \mathcal{S}_A\). Hence \(\mathcal{S}_A\) is a \(\sigma\)-algebra on \(X\).

Part (b). Suppose first that \(f\) is \(\mathcal{S}_A\)-measurable. Because \(\mathcal{S}_A \subseteq \mathcal{S}\), every set of the form \(f^{-1}(B)\) with \(B \subseteq \mathbb{R}\) Borel lies in \(\mathcal{S}_A\) and hence in \(\mathcal{S}\); thus \(f\) is \(\mathcal{S}\)-measurable. To see that \(f\) is constant on \(A\), we may assume \(A \neq \emptyset\) (otherwise there is nothing to prove). Pick \(x \in A\). The singleton \(\{f(x)\}\) is a closed, hence Borel, subset of \(\mathbb{R}\), so \(f^{-1}(\{f(x)\}) \in \mathcal{S}_A\). Since \(x \in A\) and \(x \in f^{-1}(\{f(x)\})\), we have \(A \cap f^{-1}(\{f(x)\}) \neq \emptyset\); the definition of \(\mathcal{S}_A\) therefore forces \(A \subseteq f^{-1}(\{f(x)\})\). In other words, \(f(y) = f(x)\) for every \(y \in A\), so \(f\) is constant on \(A\).

Conversely, suppose \(f\) is \(\mathcal{S}\)-measurable and \(f\) is constant on \(A\). If \(A = \emptyset\), then \(A \cap E = \emptyset\) for every \(E \in \mathcal{S}\), so \(\mathcal{S}_A = \mathcal{S}\) and there is nothing to prove. So assume \(A \neq \emptyset\) and let \(c \in \mathbb{R}\) be the value of \(f\) on \(A\). Let \(B \subseteq \mathbb{R}\) be a Borel set. Then \(f^{-1}(B) \in \mathcal{S}\) because \(f\) is \(\mathcal{S}\)-measurable. If \(c \in B\), then \(f(y) = c \in B\) for every \(y \in A\), so \(A \subseteq f^{-1}(B)\). If \(c \notin B\), then no \(y \in A\) satisfies \(f(y) \in B\), so \(A \cap f^{-1}(B) = \emptyset\). In either case \(f^{-1}(B) \in \mathcal{S}_A\). Because \(B\) was an arbitrary Borel subset of \(\mathbb{R}\), we conclude that \(f\) is \(\mathcal{S}_A\)-measurable.

Problem (2B.17)

Suppose \(X\) is a Borel subset of \(\mathbb{R}\) and \(f : X \to \mathbb{R}\) is a function such that \(\{ x \in X : f \text{ is not continuous at } x \}\) is a countable set. Prove \(f\) is a Borel measurable function.

Solution

Let

\begin{equation*} D = \{ x \in X : f \text{ is not continuous at } x \}, \end{equation*}

which by hypothesis is countable. Note first that every countable subset of \(\mathbb{R}\) is a Borel set: each singleton \(\{t\}\) is closed and hence Borel, and a countable set is a countable union of singletons, so it is Borel because the collection of Borel sets is a \(\sigma\)-algebra.

As remarked after 2.40, to prove that \(f\) is Borel measurable it suffices to show that \(f^{-1}\big((a, \infty)\big)\) is a Borel set for every \(a \in \mathbb{R}\) (this is 2.39 applied with \(\mathcal{S}\) equal to the \(\sigma\)-algebra of Borel subsets of \(\mathbb{R}\) contained in \(X\)).

So fix \(a \in \mathbb{R}\) and let

\begin{equation*} U = \{ x \in X \setminus D : f(x) > a \}, \qquad V = \{ x \in D : f(x) > a \} , \end{equation*}

so that \(f^{-1}\big((a,\infty)\big) = U \cup V\). Since \(V \subseteq D\), the set \(V\) is countable and hence Borel.

Fix \(x \in U\). Then \(f\) is continuous at \(x\) and \(f(x) > a\). Continuity of \(f\) at \(x\) means continuity relative to the domain \(X\): for every \(\varepsilon > 0\) there exists \(\delta > 0\) such that \(|f(y) - f(x)| < \varepsilon\) for all \(y \in X\) with \(|y - x| < \delta\). Applying this with \(\varepsilon = f(x) - a > 0\) produces \(\delta_x > 0\) such that \(|f(y) - f(x)| < f(x) - a\), and hence

\begin{equation*} f(y) > f(x) - \big( f(x) - a \big) = a \quad \text{for all } y \in (x - \delta_x, x + \delta_x) \cap X . \end{equation*}

Note that this conclusion applies to every \(y\) in that intersection, including points of \(D\) at which \(f\) fails to be continuous; continuity at \(x\) constrains the values of \(f\) at all nearby points of \(X\), whatever the behavior of \(f\) there. Let

\begin{equation*} G = \bigcup_{x \in U} (x - \delta_x, x + \delta_x) , \end{equation*}

which is an open subset of \(\mathbb{R}\) (a union of open intervals), hence a Borel set. Therefore \(G \cap X\) is a Borel set, being the intersection of two Borel sets.

We claim

\begin{equation*} f^{-1}\big((a,\infty)\big) = (G \cap X) \cup V . \end{equation*}

To see the inclusion \(\subseteq\): if \(x \in f^{-1}\big((a,\infty)\big)\), then either \(x \in D\), in which case \(x \in V\), or \(x \notin D\), in which case \(x \in U\) and so \(x \in (x - \delta_x, x + \delta_x) \cap X \subseteq G \cap X\). For the inclusion \(\supseteq\): certainly \(V \subseteq f^{-1}\big((a,\infty)\big)\); and if \(y \in G \cap X\), then \(y \in (x - \delta_x, x + \delta_x) \cap X\) for some \(x \in U\), and the defining property of \(\delta_x\) gives \(f(y) > a\), so \(y \in f^{-1}\big((a,\infty)\big)\). This proves the claim.

Thus \(f^{-1}\big((a,\infty)\big)\) is the union of the two Borel sets \(G \cap X\) and \(V\), and hence is a Borel set. Because \(a \in \mathbb{R}\) was arbitrary, 2.39 implies that \(f\) is a Borel measurable function.

Problem (2B.18)

Suppose \(f : \mathbb{R} \to \mathbb{R}\) is differentiable at every element of \(\mathbb{R}\). Prove that \(f^{\prime}\) is a Borel measurable function from \(\mathbb{R}\) to \(\mathbb{R}\).

Solution

Because \(f\) is differentiable at every point of \(\mathbb{R}\), it is continuous at every point of \(\mathbb{R}\). Hence \(f\) is a Borel measurable function by 2.41.

For each \(k \in \mathbb{Z}^+\) define \(g_k : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} g_k(x) = \frac{f\big(x + \tfrac{1}{k}\big) - f(x)}{1/k} = k \left( f\Big(x + \tfrac{1}{k}\Big) - f(x) \right) . \end{equation*}

The function \(x \mapsto f\big(x + \tfrac{1}{k}\big)\) is continuous on \(\mathbb{R}\) (it is the composition of the continuous translation \(x \mapsto x + \tfrac1k\) with the continuous function \(f\)), and \(f\) is continuous, so \(g_k\) is continuous on \(\mathbb{R}\). Therefore each \(g_k\) is a Borel measurable function by 2.41.

Now fix \(x \in \mathbb{R}\). By definition of the derivative,

\begin{equation*} f^{\prime}(x) = \lim_{h \to 0} \frac{f(x + h) - f(x)}{h} , \end{equation*}

and this limit exists in \(\mathbb{R}\). Taking the particular sequence \(h = \tfrac{1}{k}\), which converges to \(0\) with \(h \neq 0\), we get

\begin{equation*} \lim_{k \to \infty} g_k(x) = f^{\prime}(x) . \end{equation*}

Thus \(g_1, g_2, \ldots\) is a sequence of Borel measurable functions from \(\mathbb{R}\) to \(\mathbb{R}\) whose pointwise limit exists at every \(x \in \mathbb{R}\) and equals \(f^{\prime}(x)\).

Apply 2.48 with \(X = \mathbb{R}\) and \(\mathcal{S}\) the \(\sigma\)-algebra of Borel subsets of \(\mathbb{R}\) (for which \(\mathcal{S}\)-measurable means Borel measurable): the pointwise limit \(f^{\prime}\) of the sequence \(g_1, g_2, \ldots\) is \(\mathcal{S}\)-measurable. Hence \(f^{\prime}\) is a Borel measurable function from \(\mathbb{R}\) to \(\mathbb{R}\).

Problem (2B.19)

Suppose \(X\) is a nonempty set and \(\mathcal{S}\) is the \(\sigma\)-algebra on \(X\) consisting of all subsets of \(X\) that are either countable or have a countable complement in \(X\). Give a characterization of the \(\mathcal{S}\)-measurable real-valued functions on \(X\).

Solution

Characterization. A function \(f : X \to \mathbb{R}\) is \(\mathcal{S}\)-measurable if and only if there exists \(c \in \mathbb{R}\) such that

\begin{equation*} \{ x \in X : f(x) \neq c \} \text{ is countable} , \end{equation*}

that is, if and only if \(f\) is constant outside some countable subset of \(X\).

Proof of the easier direction. Suppose \(c \in \mathbb{R}\) and \(C = \{x \in X : f(x) \neq c\}\) is countable. Let \(B \subseteq \mathbb{R}\) be a Borel set. If \(c \in B\), then \(X \setminus C \subseteq f^{-1}(B)\), so \(X \setminus f^{-1}(B) \subseteq C\) is countable and hence \(f^{-1}(B) \in \mathcal{S}\). If \(c \notin B\), then \(f^{-1}(B) \subseteq C\), so \(f^{-1}(B)\) is countable and hence \(f^{-1}(B) \in \mathcal{S}\). In either case \(f^{-1}(B) \in \mathcal{S}\), so \(f\) is \(\mathcal{S}\)-measurable.

Proof of the other direction. Suppose \(f : X \to \mathbb{R}\) is \(\mathcal{S}\)-measurable.

If \(X\) is countable, then \(\{x \in X : f(x) \neq 0\}\) is countable and we may take \(c = 0\). So assume from now on that \(X\) is uncountable. Then no subset of \(X\) is both countable and cocountable (otherwise \(X\) would be the union of two countable sets and hence countable). Also, a countable union of countable subsets of \(X\) is countable.

For \(a \in \mathbb{R}\) write \(E_a = f^{-1}\big((a, \infty)\big) = \{x \in X : f(x) > a\}\). Each \(E_a\) belongs to \(\mathcal{S}\), so each \(E_a\) is either countable or cocountable, and (by the previous paragraph) exactly one of these holds. Let

\begin{equation*} T = \{ a \in \mathbb{R} : E_a \text{ is countable} \} . \end{equation*}

\(T\) is nonempty. Suppose instead that \(T = \emptyset\). Then \(E_n\) is cocountable for every \(n \in \mathbb{Z}^+\), so \(X \setminus \bigcap_{n=1}^{\infty} E_n = \bigcup_{n=1}^{\infty} (X \setminus E_n)\) is countable, whence \(\bigcap_{n=1}^{\infty} E_n\) is nonempty because \(X\) is uncountable. But \(\bigcap_{n=1}^{\infty} E_n = \{x \in X : f(x) > n \text{ for all } n \in \mathbb{Z}^+\} = \emptyset\), since \(f\) takes values in \(\mathbb{R}\). This contradiction shows \(T \neq \emptyset\).

\(T\) is bounded below. Because \(f\) is real-valued, \(X = \bigcup_{n=1}^{\infty} E_{-n}\). As \(X\) is uncountable and a countable union of countable sets is countable, there exists \(m \in \mathbb{Z}^+\) with \(E_{-m}\) uncountable, so \(-m \notin T\). Now observe that \(T\) is closed upward: if \(a \in T\) and \(b > a\), then \(E_b \subseteq E_a\), so \(E_b\) is countable and \(b \in T\). Consequently, if some \(a \le -m\) were in \(T\), then \(-m \in T\) (as \(-m \ge a\); the case \(-m = a\) is immediate and the case \(-m > a\) follows from upward closure), a contradiction. Hence \(a > -m\) for every \(a \in T\), so \(T\) is bounded below.

Let \(c = \inf T \in \mathbb{R}\). We claim \(\{x \in X : f(x) \neq c\}\) is countable.

First, since \(c + \tfrac1n > c = \inf T\), for each \(n \in \mathbb{Z}^+\) there is \(a \in T\) with \(a < c + \tfrac1n\), and upward closure gives \(c + \tfrac1n \in T\); that is, \(E_{c + 1/n}\) is countable. Therefore

\begin{equation*} \{ x \in X : f(x) > c \} = \bigcup_{n=1}^{\infty} E_{c + \frac{1}{n}} \end{equation*}

is a countable union of countable sets, hence countable.

Second, for each \(n \in \mathbb{Z}^+\) we have \(c - \tfrac1n < \inf T\), so \(c - \tfrac1n \notin T\); thus \(E_{c - 1/n}\) is not countable and therefore is cocountable, so \(X \setminus E_{c - 1/n} = \{x \in X : f(x) \le c - \tfrac1n\}\) is countable. Therefore

\begin{equation*} \{ x \in X : f(x) < c \} = \bigcup_{n=1}^{\infty} \Big\{ x \in X : f(x) \le c - \tfrac{1}{n} \Big\} \end{equation*}

is countable as well. (The displayed equality holds because \(f(x) < c\) if and only if \(c - f(x) > 0\), which happens if and only if \(c - f(x) \ge \tfrac1n\) for some \(n \in \mathbb{Z}^+\).)

Combining the two paragraphs, \(\{x \in X : f(x) \neq c\} = \{f > c\} \cup \{f < c\}\) is countable, as claimed. This completes the proof of the characterization.

Problem (2B.20)

Suppose \((X, \mathcal{S})\) is a measurable space and \(f, g : X \to \mathbb{R}\) are \(\mathcal{S}\)-measurable functions. Prove that if \(f(x) > 0\) for all \(x \in X\), then \(f^g\) (which is the function whose value at \(x \in X\) equals \(f(x)^{g(x)}\)) is an \(\mathcal{S}\)-measurable function.

Solution

Recall that for \(t > 0\) and \(s \in \mathbb{R}\) the power \(t^s\) is defined by \(t^s = e^{s \ln t}\). Since \(f(x) > 0\) for every \(x \in X\), this gives

\begin{equation*} f^g(x) = f(x)^{g(x)} = e^{\,g(x) \ln f(x)} \qquad \text{for all } x \in X , \end{equation*}

so \(f^g = \exp \circ \big( g \cdot (\ln \circ f) \big)\). We build this function up in three steps.

Step 1: \(\ln \circ f\) is \(\mathcal{S}\)-measurable. The set \((0, \infty)\) is an open, hence Borel, subset of \(\mathbb{R}\), and the natural logarithm \(\ln : (0, \infty) \to \mathbb{R}\) is continuous. Thus \(\ln\) is a real-valued Borel measurable function by 2.41. Because \(f(x) > 0\) for all \(x \in X\), the range of \(f\) is contained in \((0, \infty)\), which is the domain of \(\ln\). Hence 2.44 applies and shows that \(\ln \circ f : X \to \mathbb{R}\) is an \(\mathcal{S}\)-measurable function.

Step 2: \(g \cdot (\ln \circ f)\) is \(\mathcal{S}\)-measurable. Both \(g\) and \(\ln \circ f\) are \(\mathcal{S}\)-measurable real-valued functions on \(X\), so their product is \(\mathcal{S}\)-measurable by 2.46(a). Denote this product by \(h\), so \(h(x) = g(x) \ln f(x)\) for all \(x \in X\).

Step 3: \(\exp \circ h\) is \(\mathcal{S}\)-measurable. The exponential function \(\exp : \mathbb{R} \to \mathbb{R}\) is continuous on the Borel set \(\mathbb{R}\), hence is Borel measurable by 2.41. Its domain \(\mathbb{R}\) contains the range of \(h\), so 2.44 shows that \(\exp \circ h\) is an \(\mathcal{S}\)-measurable function.

Since \(f^g = \exp \circ h\) by the displayed formula above, we conclude that \(f^g\) is an \(\mathcal{S}\)-measurable function.

Problem (2B.21)

Prove 2.52. [2.52, condition for measurable function: Suppose \((X, \mathcal{S})\) is a measurable space and \(f : X \to [-\infty, \infty]\) is a function such that \(f^{-1}\big((a, \infty]\big) \in \mathcal{S}\) for all \(a \in \mathbb{R}\). Then \(f\) is an \(\mathcal{S}\)-measurable function.]

Solution

Following the pattern of the proof of 2.39, let

\begin{equation*} \mathcal{T} = \{ A \subseteq [-\infty, \infty] : f^{-1}(A) \in \mathcal{S} \} . \end{equation*}

We must show that every Borel subset of \([-\infty, \infty]\) (in the sense of 2.50) belongs to \(\mathcal{T}\).

Step 1: \(\mathcal{T}\) is a \(\sigma\)-algebra on \([-\infty, \infty]\). First, \(f^{-1}(\emptyset) = \emptyset \in \mathcal{S}\), so \(\emptyset \in \mathcal{T}\). Next, if \(A \in \mathcal{T}\), then, because \(f\) maps \(X\) into \([-\infty, \infty]\),

\begin{equation*} f^{-1}\big([-\infty, \infty] \setminus A\big) = X \setminus f^{-1}(A) \in \mathcal{S} , \end{equation*}

so \([-\infty, \infty] \setminus A \in \mathcal{T}\); thus \(\mathcal{T}\) is closed under complementation. Finally, if \(A_1, A_2, \ldots \in \mathcal{T}\), then

\begin{equation*} f^{-1}\left( \bigcup_{k=1}^{\infty} A_k \right) = \bigcup_{k=1}^{\infty} f^{-1}(A_k) \in \mathcal{S} , \end{equation*}

so \(\bigcup_{k=1}^{\infty} A_k \in \mathcal{T}\); thus \(\mathcal{T}\) is closed under countable unions. Hence \(\mathcal{T}\) is a \(\sigma\)-algebra on \([-\infty, \infty]\). In particular, \(\mathcal{T}\) is closed under countable intersections and under set differences (by 2.25).

Step 2: \(\{\infty\}\), \(\{-\infty\}\), and \(\mathbb{R}\) belong to \(\mathcal{T}\). By hypothesis \((a, \infty] \in \mathcal{T}\) for every \(a \in \mathbb{R}\). Now

\begin{equation*} \bigcap_{n=1}^{\infty} (n, \infty] = \{\infty\} , \end{equation*}

because an element \(t \in [-\infty, \infty]\) satisfies \(t > n\) for every \(n \in \mathbb{Z}^+\) if and only if \(t = \infty\). Hence \(\{\infty\} \in \mathcal{T}\). Similarly

\begin{equation*} \bigcup_{n=1}^{\infty} (-n, \infty] = (-\infty, \infty] , \end{equation*}

because every \(t \in [-\infty, \infty]\) other than \(-\infty\) satisfies \(t > -n\) for some \(n \in \mathbb{Z}^+\). Hence \((-\infty, \infty] \in \mathcal{T}\), and taking complements gives \(\{-\infty\} \in \mathcal{T}\). Therefore

\begin{equation*} \mathbb{R} = [-\infty, \infty] \setminus \big( \{\infty\} \cup \{-\infty\} \big) \in \mathcal{T} . \end{equation*}

Step 3: every Borel subset of \(\mathbb{R}\) belongs to \(\mathcal{T}\). Let

\begin{equation*} \mathcal{T}_{\mathbb{R}} = \{ A \in \mathcal{T} : A \subseteq \mathbb{R} \} . \end{equation*}

We check that \(\mathcal{T}_{\mathbb{R}}\) is a \(\sigma\)-algebra on \(\mathbb{R}\). Indeed \(\emptyset \in \mathcal{T}_{\mathbb{R}}\); if \(A \in \mathcal{T}_{\mathbb{R}}\), then

\begin{equation*} \mathbb{R} \setminus A = \mathbb{R} \cap \big( [-\infty, \infty] \setminus A \big) \in \mathcal{T} \end{equation*}

because \(\mathbb{R} \in \mathcal{T}\) by Step 2 and \(\mathcal{T}\) is closed under complementation and finite intersections, and clearly \(\mathbb{R} \setminus A \subseteq \mathbb{R}\), so \(\mathbb{R} \setminus A \in \mathcal{T}_{\mathbb{R}}\); and \(\mathcal{T}_{\mathbb{R}}\) is closed under countable unions because \(\mathcal{T}\) is and because a countable union of subsets of \(\mathbb{R}\) is a subset of \(\mathbb{R}\).

Moreover, for every \(a \in \mathbb{R}\) we have

\begin{equation*} (a, \infty) = (a, \infty] \setminus \{\infty\} \in \mathcal{T} , \end{equation*}

and \((a, \infty) \subseteq \mathbb{R}\), so \((a, \infty) \in \mathcal{T}_{\mathbb{R}}\).

Thus \(\mathcal{T}_{\mathbb{R}}\) is a \(\sigma\)-algebra on \(\mathbb{R}\) containing \(\{ (a, \infty) : a \in \mathbb{R} \}\). The second half of the proof of 2.39 shows that any such \(\sigma\)-algebra contains every Borel subset of \(\mathbb{R}\); for completeness we repeat the argument. Closure under complementation gives \((-\infty, b] \in \mathcal{T}_{\mathbb{R}}\) for all \(b \in \mathbb{R}\), and then closure under finite intersections (2.25) gives \((a, b] = (a, \infty) \cap (-\infty, b] \in \mathcal{T}_{\mathbb{R}}\) for all \(a, b \in \mathbb{R}\). Because

\begin{equation*} (a, b) = \bigcup_{k=1}^{\infty} \Big( a, b - \tfrac{1}{k} \Big] \quad \text{and} \quad (-\infty, b) = \bigcup_{k=1}^{\infty} \Big( -k, b - \tfrac{1}{k} \Big] , \end{equation*}

and because \((a, \infty)\) and \(\mathbb{R}\) are already in \(\mathcal{T}_{\mathbb{R}}\), the \(\sigma\)-algebra \(\mathcal{T}_{\mathbb{R}}\) contains every open interval of \(\mathbb{R}\) (bounded or unbounded). Every open subset of \(\mathbb{R}\) is a countable union of open intervals: if \(G \subseteq \mathbb{R}\) is open, then for each \(x \in G\) there are rationals \(p < x < q\) with \((p, q) \subseteq G\), and \(G\) is the union of the countably many such rational intervals. Hence \(\mathcal{T}_{\mathbb{R}}\) contains every open subset of \(\mathbb{R}\), and therefore, being a \(\sigma\)-algebra, it contains the smallest \(\sigma\)-algebra on \(\mathbb{R}\) containing the open sets, which is the collection of Borel subsets of \(\mathbb{R}\).

Step 4: conclusion. Let \(C \subseteq [-\infty, \infty]\) be a Borel set. By 2.50, \(B := C \cap \mathbb{R}\) is a Borel subset of \(\mathbb{R}\), and \(C\) equals one of

\begin{equation*} B, \quad B \cup \{\infty\}, \quad B \cup \{-\infty\}, \quad B \cup \{\infty, -\infty\} . \end{equation*}

By Step 3 we have \(B \in \mathcal{T}_{\mathbb{R}} \subseteq \mathcal{T}\), by Step 2 we have \(\{\infty\}, \{-\infty\} \in \mathcal{T}\), and \(\mathcal{T}\) is closed under finite unions. Hence \(C \in \mathcal{T}\) in all four cases, which means \(f^{-1}( C) \in \mathcal{S}\).

Because \(C\) was an arbitrary Borel subset of \([-\infty, \infty]\), the function \(f\) is \(\mathcal{S}\)-measurable, completing the proof of 2.52.

Problem (2B.22)

Suppose \(B \subseteq \mathbb{R}\) and \(f : B \to \mathbb{R}\) is an increasing function. Prove that \(f\) is continuous at every element of \(B\) except for a countable subset of \(B\).

Solution

Recall that \(f\) increasing means \(f(s) \le f(t)\) for all \(s, t \in B\) with \(s < t\), and that continuity of \(f\) at \(x \in B\) means: for every \(\varepsilon > 0\) there exists \(\delta > 0\) such that \(|f(t) - f(x)| < \varepsilon\) for all \(t \in B \cap (x - \delta, x + \delta)\) (continuity is relative to the domain \(B\)).

Call \(x \in B\) a left limit point of \(B\) if \(B \cap (x - \delta, x) \neq \emptyset\) for every \(\delta > 0\), and a right limit point of \(B\) if \(B \cap (x, x + \delta) \neq \emptyset\) for every \(\delta > 0\).

For \(x \in B\) define

\begin{equation*} \alpha(x) = \begin{cases} \sup\{f(t) : t \in B,\ t < x\} & \text{if } B \cap (-\infty, x) \neq \emptyset, \\ f(x) & \text{otherwise,}\end{cases} \qquad \beta(x) = \begin{cases} \inf\{f(t) : t \in B,\ t > x\} & \text{if } B \cap (x, \infty) \neq \emptyset, \\ f(x) & \text{otherwise.}\end{cases} \end{equation*}

Because \(f\) is increasing, \(f(t) \le f(x)\) for every \(t \in B\) with \(t < x\), so the set defining \(\alpha(x)\) is a nonempty set of real numbers bounded above by \(f(x)\); hence \(\alpha(x) \in \mathbb{R}\) and \(\alpha(x) \le f(x)\). Symmetrically \(\beta(x) \in \mathbb{R}\) and \(f(x) \le \beta(x)\).

Step 1: a characterization of the discontinuity points. Let \(D\) denote the set of \(x \in B\) at which \(f\) is not continuous. We claim that \(x \in D\) if and only if at least one of the following holds:

  • (i) \(x\) is a left limit point of \(B\) and \(\alpha(x) < f(x)\);
  • (ii) \(x\) is a right limit point of \(B\) and \(f(x) < \beta(x)\).

First suppose neither (i) nor (ii) holds; we show \(f\) is continuous at \(x\). Let \(\varepsilon > 0\). We produce \(\delta_1 > 0\) controlling the left side. If \(x\) is not a left limit point of \(B\), choose \(\delta_1 > 0\) with \(B \cap (x - \delta_1, x) = \emptyset\); then the only element of \(B \cap (x - \delta_1, x]\) is \(x\) itself, so \(|f(t) - f(x)| = 0 < \varepsilon\) for every \(t \in B \cap (x - \delta_1, x]\). If \(x\) is a left limit point of \(B\), then \(B \cap (-\infty, x) \neq \emptyset\), so \(\alpha(x)\) is genuinely the supremum written above; the failure of (i), combined with \(\alpha(x) \le f(x)\), forces \(\alpha(x) = f(x)\). By the definition of supremum there exists \(t_0 \in B\) with \(t_0 < x\) and \(f(t_0) > f(x) - \varepsilon\). Put \(\delta_1 = x - t_0 > 0\). For \(t \in B \cap (x - \delta_1, x]\) we have \(t_0 < t \le x\), hence \(f(x) - \varepsilon < f(t_0) \le f(t) \le f(x)\), so \(|f(t) - f(x)| < \varepsilon\).

Symmetrically, using that (ii) fails, there exists \(\delta_2 > 0\) such that \(|f(t) - f(x)| < \varepsilon\) for all \(t \in B \cap [x, x + \delta_2)\). Taking \(\delta = \min\{\delta_1, \delta_2\}\) gives \(|f(t) - f(x)| < \varepsilon\) for all \(t \in B \cap (x - \delta, x + \delta)\). Thus \(f\) is continuous at \(x\), so \(x \notin D\).

Conversely, suppose (i) holds and let \(\varepsilon = f(x) - \alpha(x) > 0\). Given any \(\delta > 0\), since \(x\) is a left limit point there exists \(t \in B \cap (x - \delta, x)\), and then \(f(t) \le \alpha(x) = f(x) - \varepsilon\), so \(|f(t) - f(x)| \ge \varepsilon\). Hence \(f\) is not continuous at \(x\). The case where (ii) holds is symmetric. This proves the claim.

Step 2: assigning disjoint intervals. For \(x \in D\) define a nonempty open interval \(I_x \subseteq \mathbb{R}\) by

\begin{equation*} I_x = \begin{cases} \bigl(\alpha(x),\, f(x)\bigr) & \text{if (i) holds for } x, \\ \bigl(f(x),\, \beta(x)\bigr) & \text{if (i) fails for } x \text{ (so (ii) holds).}\end{cases} \end{equation*}

In either case \(I_x\) is a nonempty bounded open interval.

We claim that \(I_x \cap I_y = \emptyset\) whenever \(x, y \in D\) with \(x < y\). It suffices to show \(\sup I_x \le \inf I_y\).

Case A: \(I_x = (\alpha(x), f(x))\), so \(\sup I_x = f(x)\). If \(I_y = (\alpha(y), f(y))\), then since \(x \in B\) and \(x < y\) we get \(\inf I_y = \alpha(y) = \sup\{f(t) : t \in B,\ t < y\} \ge f(x) = \sup I_x\). If instead \(I_y = (f(y), \beta(y))\), then \(\inf I_y = f(y) \ge f(x) = \sup I_x\) because \(f\) is increasing.

Case B: \(I_x = (f(x), \beta(x))\), so (ii) holds for \(x\) and \(\sup I_x = \beta(x)\). If \(I_y = (f(y), \beta(y))\), then \(\inf I_y = f(y) \ge \beta(x)\), because \(y \in B\) and \(y > x\) imply \(\beta(x) = \inf\{f(t) : t \in B,\ t > x\} \le f(y)\). If instead \(I_y = (\alpha(y), f(y))\), use that \(x\) is a right limit point of \(B\): taking \(\delta = y - x\) gives some \(s \in B\) with \(x < s < y\). Then \(\beta(x) \le f(s)\) (as \(s > x\)) and \(f(s) \le \alpha(y)\) (as \(s < y\)), so \(\sup I_x = \beta(x) \le \alpha(y) = \inf I_y\).

Thus \(\{I_x : x \in D\}\) is a pairwise disjoint collection of nonempty open intervals, and \(x \mapsto I_x\) is injective on \(D\) (distinct \(x\) give disjoint, hence distinct, nonempty intervals).

Step 3: countability. For each \(x \in D\) choose (using the density of \(\mathbb{Q}\) in \(\mathbb{R}\)) a rational number \(q_x \in I_x\). If \(x, y \in D\) with \(x \neq y\), then \(I_x \cap I_y = \emptyset\), so \(q_x \neq q_y\). Hence \(x \mapsto q_x\) is an injective map from \(D\) into \(\mathbb{Q}\), and therefore \(D\) is countable.

So \(f\) is continuous at every element of \(B\) except for the countable set \(D \subseteq B\).

Problem (2B.23)

Suppose \(f : \mathbb{R} \to \mathbb{R}\) is a strictly increasing function. Prove that the inverse function \(f^{-1} : f(\mathbb{R}) \to \mathbb{R}\) is a continuous function.

[Note that this exercise does not have as a hypothesis that \(f\) is continuous.]

Solution

First note that \(f\) is injective: if \(x \neq y\), say \(x < y\), then \(f(x) < f(y)\), so \(f(x) \neq f(y)\). Hence \(f\) is a bijection from \(\mathbb{R}\) onto \(f(\mathbb{R})\) and the inverse function \(g = f^{-1} : f(\mathbb{R}) \to \mathbb{R}\) is well defined; it is characterized by \(g(f(x)) = x\) for all \(x \in \mathbb{R}\).

We also record the following consequence of strict monotonicity: for \(x, u \in \mathbb{R}\),

\begin{equation*} f(x) < f(u) \implies x < u. \end{equation*}

Indeed, if \(x \ge u\) then \(f(x) \ge f(u)\) (with equality only when \(x = u\)), contradicting \(f(x) < f(u)\).

Now fix \(y_0 \in f(\mathbb{R})\) and let \(x_0 = g(y_0)\), so \(y_0 = f(x_0)\). Let \(\varepsilon > 0\). Because \(f\) is strictly increasing,

\begin{equation*} f(x_0 - \varepsilon) < f(x_0) = y_0 < f(x_0 + \varepsilon). \end{equation*}

Hence

\begin{equation*} \delta = \min\bigl\{\, y_0 - f(x_0 - \varepsilon),\ f(x_0 + \varepsilon) - y_0 \,\bigr\} > 0 . \end{equation*}

Suppose \(y \in f(\mathbb{R})\) with \(|y - y_0| < \delta\). Then

\begin{equation*} f(x_0 - \varepsilon) \le y_0 - \delta < y < y_0 + \delta \le f(x_0 + \varepsilon). \end{equation*}

Write \(x = g(y)\), so \(y = f(x)\). From \(f(x_0 - \varepsilon) < f(x)\) we get \(x_0 - \varepsilon < x\) by the displayed implication above, and from \(f(x) < f(x_0 + \varepsilon)\) we get \(x < x_0 + \varepsilon\). Therefore

\begin{equation*} |g(y) - g(y_0)| = |x - x_0| < \varepsilon . \end{equation*}

Thus for every \(y_0 \in f(\mathbb{R})\) and every \(\varepsilon > 0\) there exists \(\delta > 0\) such that \(|g(y) - g(y_0)| < \varepsilon\) for all \(y \in f(\mathbb{R})\) with \(|y - y_0| < \delta\). In other words, \(g = f^{-1}\) is continuous at every point of its domain \(f(\mathbb{R})\), as desired.

Note that continuity of \(f\) was never used; the only inputs were that \(f\) is strictly increasing and that its domain is all of \(\mathbb{R}\) (the latter guaranteeing that \(x_0 \pm \varepsilon\) lie in the domain).

Problem (2B.24)

Suppose \(f : \mathbb{R} \to \mathbb{R}\) is a strictly increasing function and \(B \subseteq \mathbb{R}\) is a Borel set. Prove that \(f(B)\) is a Borel set.

Solution

Let \(g = f^{-1} : f(\mathbb{R}) \to \mathbb{R}\) be the inverse function, which exists because \(f\) is injective (strictly increasing), and which is continuous by Exercise 23 of this section. The key observation is

\begin{equation*} g^{-1}(B) = \{\, y \in f(\mathbb{R}) : g(y) \in B \,\} = \{\, f(x) : x \in \mathbb{R},\ x \in B \,\} = f(B), \end{equation*}

where the middle equality holds because every \(y \in f(\mathbb{R})\) has the form \(y = f(x)\) for a unique \(x \in \mathbb{R}\), and then \(g(y) = x\).

So by 2.41 (every continuous real-valued function defined on a Borel subset of \(\mathbb{R}\) is Borel measurable) it suffices to prove that the domain \(f(\mathbb{R})\) of \(g\) is a Borel set: then \(g\) is Borel measurable, hence \(f(B) = g^{-1}(B)\) is a Borel set for every Borel set \(B \subseteq \mathbb{R}\).

Proof that \(f(\mathbb{R})\) is a Borel set. For \(x \in \mathbb{R}\) put

\begin{equation*} f(x^-) = \sup\{f(t) : t < x\}, \qquad f(x^+) = \inf\{f(t) : t > x\}, \end{equation*}

which are real numbers with \(f(x^-) \le f(x) \le f(x^+)\) (the sets involved are nonempty and bounded by \(f(x)\) on the appropriate side). Put

\begin{equation*} L = \inf\{f(t) : t \in \mathbb{R}\} \in [-\infty, \infty), \qquad M = \sup\{f(t) : t \in \mathbb{R}\} \in (-\infty, \infty], \end{equation*}

with the convention that \((-\infty, L] = \emptyset\) if \(L = -\infty\) and \([M, \infty) = \emptyset\) if \(M = \infty\). Finally let

\begin{equation*} D = \{\, x \in \mathbb{R} : f(x^-) < f(x) \text{ or } f(x) < f(x^+) \,\}, \end{equation*}

which is countable. Indeed, every real number is both a left limit point and a right limit point of \(\mathbb{R}\), so the characterization proved in Step 1 of Exercise 22 of this section (applied with \(B = \mathbb{R}\), where \(\alpha(x) = f(x^-)\) and \(\beta(x) = f(x^+)\)) says that \(D\) is exactly the set of points at which \(f\) is not continuous; that exercise shows this set is countable.

We claim

\begin{equation*} \mathbb{R} \setminus f(\mathbb{R}) = (-\infty, L] \ \cup\ [M, \infty) \ \cup \bigcup_{x \in D} \Bigl( \bigl[f(x^-), f(x)\bigr) \cup \bigl(f(x), f(x^+)\bigr] \Bigr). \end{equation*}

First, none of the sets on the right meets \(f(\mathbb{R})\).

  • If \(L \in \mathbb{R}\) and \(f(u) = L\) for some \(u\), then \(f(u - 1) < f(u) = L\), contradicting the definition of \(L\) as an infimum. So the infimum \(L\) is not attained, and since \(f(t) \ge L\) for all \(t\), we get \((-\infty, L] \cap f(\mathbb{R}) = \emptyset\). Symmetrically \([M, \infty) \cap f(\mathbb{R}) = \emptyset\).
  • Fix \(x\) and suppose \(f(x) < f(x^+)\). The infimum \(f(x^+)\) is not attained: if \(f(u) = f(x^+)\) with \(u > x\), then for \(t\) with \(x < t < u\) we would have \(f(t) < f(u) = f(x^+) = \inf\{f(s) : s > x\}\), a contradiction; and if \(u \le x\) then \(f(u) \le f(x) < f(x^+)\). Moreover any value \(f(u)\) lying in \((f(x), f(x^+)]\) must have \(u > x\) (since \(u \le x\) gives \(f(u) \le f(x)\)), whence \(f(u) \ge f(x^+)\), so \(f(u) = f(x^+)\), which we just excluded. Hence \((f(x), f(x^+)] \cap f(\mathbb{R}) = \emptyset\). Symmetrically \([f(x^-), f(x)) \cap f(\mathbb{R}) = \emptyset\).

For the reverse inclusion, let \(y \notin f(\mathbb{R})\) and suppose \(L < y < M\) (otherwise \(y \le L\) or \(y \ge M\), so \(y\) lies in one of the first two sets). Let

\begin{equation*} T = \{\, t \in \mathbb{R} : f(t) < y \,\}. \end{equation*}

Since \(y > L\), the set \(T\) is nonempty. Since \(y < M\), there is \(s\) with \(f(s) > y\); for \(t \in T\) we then have \(f(t) < y < f(s)\), which forces \(t < s\) (as \(t \ge s\) would give \(f(t) \ge f(s)\)). So \(T\) is bounded above, and \(x = \sup T \in \mathbb{R}\).

If \(t < x\), then there is \(t^{\prime} \in T\) with \(t < t^{\prime} \le x\), so \(f(t) \le f(t^{\prime}) < y\); taking the supremum over \(t < x\) gives \(f(x^-) \le y\). If \(t > x\), then \(t \notin T\), so \(f(t) \ge y\), and in fact \(f(t) > y\) because \(y \notin f(\mathbb{R})\); taking the infimum over \(t > x\) gives \(f(x^+) \ge y\). Also \(f(x) \neq y\) since \(y \notin f(\mathbb{R})\). There are two cases.

  • If \(f(x) > y\), then \(f(x^-) \le y < f(x)\), so \(y \in [f(x^-), f(x))\) and in particular \(f(x^-) < f(x)\), so \(x \in D\).
  • If \(f(x) < y\), then \(f(x) < y \le f(x^+)\), so \(y \in (f(x), f(x^+)]\) and in particular \(f(x) < f(x^+)\), so \(x \in D\).

Either way \(y\) belongs to the right side, proving the claimed equality.

The right side is a union of \(2\) intervals together with a countable family of intervals indexed by the countable set \(D\); each interval is a Borel set, and the collection of Borel sets is closed under countable unions (it is a \(\sigma\)-algebra), so \(\mathbb{R} \setminus f(\mathbb{R})\) is a Borel set. Hence \(f(\mathbb{R})\) is a Borel set, being the complement of a Borel set.

This completes the proof: \(g\) is a continuous function on the Borel set \(f(\mathbb{R})\), hence Borel measurable by 2.41, hence \(f(B) = g^{-1}(B)\) is a Borel set.

Problem (2B.25)

Suppose \(B \subseteq \mathbb{R}\) and \(f : B \to \mathbb{R}\) is an increasing function. Prove that there exists a sequence \(f_1, f_2, \ldots\) of strictly increasing functions from \(B\) to \(\mathbb{R}\) such that

\begin{equation*} f(x) = \lim_{k \to \infty} f_k(x) \end{equation*}

for every \(x \in B\).

Solution

For each \(k \in \mathbb{Z}^+\) define \(f_k : B \to \mathbb{R}\) by

\begin{equation*} f_k(x) = f(x) + \frac{x}{k}. \end{equation*}

Each \(f_k\) is a well-defined real-valued function on \(B\).

\(f_k\) is strictly increasing: suppose \(x, y \in B\) with \(x < y\). Because \(f\) is increasing, \(f(x) \le f(y)\). Because \(k > 0\) and \(x < y\), we have \(x/k < y/k\). Adding the two inequalities gives

\begin{equation*} f_k(x) = f(x) + \frac{x}{k} < f(y) + \frac{y}{k} = f_k(y). \end{equation*}

Thus \(f_k(x) < f_k(y)\) whenever \(x < y\), i.e. \(f_k\) is strictly increasing.

Pointwise convergence: fix \(x \in B\). Then \(x\) is a fixed real number, so

\begin{equation*} \lim_{k \to \infty} f_k(x) = \lim_{k \to \infty} \Bigl( f(x) + \frac{x}{k} \Bigr) = f(x) + \lim_{k \to \infty} \frac{x}{k} = f(x) + 0 = f(x). \end{equation*}

Hence \(f_1, f_2, \ldots\) is a sequence of strictly increasing functions from \(B\) to \(\mathbb{R}\) converging pointwise on \(B\) to \(f\), as required.

Remark. Any strictly increasing function \(\varphi : \mathbb{R} \to \mathbb{R}\) would serve equally well in place of the identity: setting \(f_k(x) = f(x) + \varphi(x)/k\) gives strictly increasing functions converging pointwise to \(f\). Taking \(\varphi\) bounded, say \(\varphi(x) = \arctan x\), has the extra feature that the convergence is uniform on \(B\) and that each \(f_k\) is bounded whenever \(f\) is.

Problem (2B.26)

Suppose \(B \subseteq \mathbb{R}\) and \(f : B \to \mathbb{R}\) is a bounded increasing function. Prove that there exists an increasing function \(g : \mathbb{R} \to \mathbb{R}\) such that \(g(x) = f(x)\) for all \(x \in B\).

Solution

If \(B = \emptyset\) then \(f\) is the empty function and the constant function \(g = 0\) on \(\mathbb{R}\) works (it is increasing and the condition on \(B\) is vacuous). So assume \(B \neq \emptyset\).

Because \(f\) is bounded, the set \(f(B) = \{f(t) : t \in B\}\) is a nonempty bounded set of real numbers; let

\begin{equation*} m = \inf f(B) \in \mathbb{R}, \qquad M = \sup f(B) \in \mathbb{R}. \end{equation*}

Define \(g : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} g(x) = \begin{cases} m & \text{if } B \cap (-\infty, x] = \emptyset, \\[2pt] \sup\{\, f(t) : t \in B,\ t \le x \,\} & \text{if } B \cap (-\infty, x] \neq \emptyset. \end{cases} \end{equation*}

\(g\) is well defined and real-valued: in the second case the set \(\{f(t) : t \in B,\ t \le x\}\) is a nonempty subset of \(f(B)\), hence bounded above by \(M\), so its supremum is a real number lying in \([m, M]\). In the first case \(g(x) = m \in \mathbb{R}\). In particular \(m \le g(x) \le M\) for all \(x \in \mathbb{R}\).

\(g\) is increasing: suppose \(x < y\). If \(B \cap (-\infty, x] = \emptyset\), then \(g(x) = m \le g(y)\), because \(g(y)\) is either \(m\) itself or the supremum of a nonempty subset of \(f(B)\), and every element of \(f(B)\) is at least \(m\). If \(B \cap (-\infty, x] \neq \emptyset\), then \(B \cap (-\infty, y] \neq \emptyset\) as well and

\begin{equation*} \{\, f(t) : t \in B,\ t \le x \,\} \subseteq \{\, f(t) : t \in B,\ t \le y \,\}, \end{equation*}

so the supremum over the larger set is at least the supremum over the smaller set, i.e. \(g(x) \le g(y)\).

\(g\) extends \(f\): let \(x \in B\). Then \(x \in B \cap (-\infty, x]\), so the second case of the definition applies and

\begin{equation*} g(x) = \sup\{\, f(t) : t \in B,\ t \le x \,\}. \end{equation*}

For \(t \in B\) with \(t \le x\) we have \(f(t) \le f(x)\): this is clear if \(t = x\), and follows from \(f\) being increasing if \(t < x\). So \(f(x)\) is an upper bound of the set, and \(f(x)\) belongs to the set (take \(t = x\)). Hence the supremum equals \(f(x)\), i.e. \(g(x) = f(x)\).

Thus \(g : \mathbb{R} \to \mathbb{R}\) is an increasing function with \(g|_B = f\), as required.

Remark. Boundedness is used only to guarantee that \(g\) takes values in \(\mathbb{R}\) rather than in \([-\infty, \infty]\); without it the same formula still defines an increasing function from \(\mathbb{R}\) to \([-\infty, \infty]\) extending \(f\). Some hypothesis of this kind is genuinely needed: let \(B = (0,1)\) and \(f(x) = 1/(1 - x)\), which is increasing on \(B\) but unbounded above. If \(g : \mathbb{R} \to \mathbb{R}\) were an increasing extension, then \(g(1) \ge g(x) = 1/(1 - x)\) for every \(x \in (0,1)\), and letting \(x\) increase to \(1\) makes the right side arbitrarily large, so no real number \(g(1)\) can work.

Problem (2B.27)

Prove or give a counterexample: If \((X, \mathcal{S})\) is a measurable space and

\begin{equation*} f : X \to [-\infty, \infty] \end{equation*}

is a function such that \(f^{-1}\bigl((a, \infty)\bigr) \in \mathcal{S}\) for every \(a \in \mathbb{R}\), then \(f\) is an \(\mathcal{S}\)-measurable function.

Solution

The statement is false; here is a counterexample.

Let \(X = \{1, 2\}\) and \(\mathcal{S} = \{\emptyset, X\}\), which is a \(\sigma\)-algebra on \(X\). Define \(f : X \to [-\infty, \infty]\) by

\begin{equation*} f(1) = \infty, \qquad f(2) = -\infty . \end{equation*}

The hypothesis holds. Indeed, for \(a \in \mathbb{R}\),

\begin{equation*} f^{-1}\bigl((a, \infty)\bigr) = \{\, x \in X : a < f(x) < \infty \,\} = \emptyset \in \mathcal{S}, \end{equation*}

because neither value of \(f\) lies in \(\mathbb{R}\): \(f(1) = \infty \notin (a, \infty)\) and \(f(2) = -\infty \notin (a, \infty)\).

The conclusion fails. By 2.50, a subset of \([-\infty, \infty]\) is a Borel set exactly when its intersection with \(\mathbb{R}\) is a Borel subset of \(\mathbb{R}\); since \(\{\infty\} \cap \mathbb{R} = \emptyset\) is a Borel set, \(\{\infty\}\) is a Borel subset of \([-\infty, \infty]\). But

\begin{equation*} f^{-1}\bigl(\{\infty\}\bigr) = \{1\} \notin \mathcal{S} . \end{equation*}

Hence \(f\) is not \(\mathcal{S}\)-measurable in the sense of 2.51.

What goes wrong, and what remains true. Let

\begin{equation*} \mathcal{T} = \{\, A \subseteq [-\infty, \infty] : f^{-1}(A) \in \mathcal{S} \,\}. \end{equation*}

Exactly as in the proof of 2.39 (using 2.33), \(\mathcal{T}\) is a \(\sigma\)-algebra on \([-\infty, \infty]\). The hypothesis says \((a, \infty) \in \mathcal{T}\) for all \(a \in \mathbb{R}\). For \(a < b\) the complement of \((b, \infty)\) in \([-\infty, \infty]\) is \([-\infty, b] \cup \{\infty\}\), and intersecting it with \((a, \infty)\) gives

\begin{equation*} (a, \infty) \cap \bigl( [-\infty, b] \cup \{\infty\} \bigr) = (a, b], \end{equation*}

since \(\infty \notin (a, \infty)\). Hence \((a, b] \in \mathcal{T}\) for all \(a < b\), and therefore \(\mathbb{R} = \bigcup_{k=1}^{\infty} (-k, k] \in \mathcal{T}\). It follows that

\begin{equation*} \mathcal{T}|_{\mathbb{R}} = \{\, A \subseteq \mathbb{R} : A \in \mathcal{T} \,\} \end{equation*}

is a \(\sigma\)-algebra on \(\mathbb{R}\): it contains \(\emptyset\), it is closed under countable unions because \(\mathcal{T}\) is, and if \(A \in \mathcal{T}|_{\mathbb{R}}\) then \(\mathbb{R} \setminus A = \bigl( [-\infty, \infty] \setminus A \bigr) \cap \mathbb{R} \in \mathcal{T}|_{\mathbb{R}}\) using \(\mathbb{R} \in \mathcal{T}\) and 2.25(b). Every open subset of \(\mathbb{R}\) is a countable union of sets of the form \((a, b]\) (write \((c, d) = \bigcup_{k} (c, d_k]\) with \(d_k\) increasing to \(d\), and use that every open subset of \(\mathbb{R}\) is a countable union of open intervals), so \(\mathcal{T}|_{\mathbb{R}}\) contains every open subset of \(\mathbb{R}\) and hence every Borel subset of \(\mathbb{R}\). So the hypothesis does imply

\begin{equation*} f^{-1}(B) \in \mathcal{S} \quad \text{for every Borel set } B \subseteq \mathbb{R}. \end{equation*}

What it does not control is the two points \(\infty\) and \(-\infty\): no countable combination of the sets \((a, \infty)\), \(a \in \mathbb{R}\), isolates \(\{\infty\}\), and the counterexample above exploits precisely this. Replacing the hypothesis by \(f^{-1}\bigl((a, \infty]\bigr) \in \mathcal{S}\) for all \(a \in \mathbb{R}\) repairs the statement: that is 2.52, since then

\begin{equation*} f^{-1}\bigl(\{\infty\}\bigr) = \bigcap_{k=1}^{\infty} f^{-1}\bigl((k, \infty]\bigr) \in \mathcal{S} \end{equation*}

and \(f^{-1}\bigl(\{-\infty\}\bigr) = X \setminus \bigcup_{k=1}^{\infty} f^{-1}\bigl((-k, \infty]\bigr) \in \mathcal{S}\).

Problem (2B.28)

Suppose \(f : B \to \mathbb{R}\) is a Borel measurable function. Define \(g : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} g(x) = \begin{cases} f(x) & \text{if } x \in B, \\ 0 & \text{if } x \in \mathbb{R} \setminus B. \end{cases} \end{equation*}

Prove that \(g\) is a Borel measurable function.

Solution

Here \(B \subseteq \mathbb{R}\) and, by 2.40, \(f\) being Borel measurable means that \(f^{-1}(A)\) is a Borel set for every Borel set \(A \subseteq \mathbb{R}\).

First note that \(B\) itself is a Borel set: as observed after 2.40,

\begin{equation*} B = f^{-1}(\mathbb{R}), \end{equation*}

and \(\mathbb{R}\) is a Borel subset of \(\mathbb{R}\), so \(B\) is a Borel set. Consequently \(\mathbb{R} \setminus B\) is a Borel set, since the collection of Borel sets is a \(\sigma\)-algebra on \(\mathbb{R}\) and hence closed under complementation.

Now let \(A \subseteq \mathbb{R}\) be a Borel set. We compute \(g^{-1}(A)\) by splitting \(\mathbb{R}\) into \(B\) and \(\mathbb{R} \setminus B\):

\begin{equation*} g^{-1}(A) = \bigl( g^{-1}(A) \cap B \bigr) \cup \bigl( g^{-1}(A) \cap (\mathbb{R} \setminus B) \bigr). \end{equation*}

For \(x \in B\) we have \(g(x) = f(x)\), so \(g^{-1}(A) \cap B = f^{-1}(A)\) (note \(f^{-1}(A) \subseteq B\) automatically, since \(B\) is the domain of \(f\)). For \(x \in \mathbb{R} \setminus B\) we have \(g(x) = 0\), so

\begin{equation*} g^{-1}(A) \cap (\mathbb{R} \setminus B) = \begin{cases} \mathbb{R} \setminus B & \text{if } 0 \in A, \\ \emptyset & \text{if } 0 \notin A. \end{cases} \end{equation*}

Therefore

\begin{equation*} g^{-1}(A) = \begin{cases} f^{-1}(A) \cup (\mathbb{R} \setminus B) & \text{if } 0 \in A, \\ f^{-1}(A) & \text{if } 0 \notin A. \end{cases} \end{equation*}

In both cases \(g^{-1}(A)\) is a Borel set: \(f^{-1}(A)\) is a Borel set because \(f\) is Borel measurable, \(\mathbb{R} \setminus B\) is a Borel set as shown above, and the union of two Borel sets is a Borel set by 2.25(b).

Since \(g^{-1}(A)\) is a Borel set for every Borel set \(A \subseteq \mathbb{R}\), the function \(g : \mathbb{R} \to \mathbb{R}\) is Borel measurable, as desired.

(Alternatively, one could verify only the sets \((a, \infty)\), as permitted by 2.39 and the remark following 2.40: for \(a \ge 0\) we get \(g^{-1}\bigl((a, \infty)\bigr) = f^{-1}\bigl((a, \infty)\bigr)\), and for \(a < 0\) we get \(g^{-1}\bigl((a, \infty)\bigr) = f^{-1}\bigl((a, \infty)\bigr) \cup (\mathbb{R} \setminus B)\), in each case a Borel set.)

Problem (2B.29)

Give an example of a measurable space \((X, \mathcal{S})\) and a family \(\{f_t\}_{t \in \mathbb{R}}\) such that each \(f_t\) is an \(\mathcal{S}\)-measurable function from \(X\) to \([0, 1]\), but the function \(f : X \to [0, 1]\) defined by

\begin{equation*} f(x) = \sup\{f_t(x) : t \in \mathbb{R}\} \end{equation*}

is not \(\mathcal{S}\)-measurable.

[Compare this exercise to 2.53, where the index set is \(\mathbb{Z}^{+}\) rather than \(\mathbb{R}\).]

Solution

Take \(X = \mathbb{R}\) and let

\begin{equation*} \mathcal{S} = \{E \subseteq \mathbb{R} : E \text{ is countable or } \mathbb{R} \setminus E \text{ is countable}\}. \end{equation*}

By the first bullet point of 2.28, \(\mathcal{S}\) is a \(\sigma\)-algebra on \(\mathbb{R}\) (indeed it is the smallest \(\sigma\)-algebra on \(\mathbb{R}\) containing all the singletons), so \((X, \mathcal{S})\) is a measurable space.

For \(t \in \mathbb{R}\) define \(f_t : \mathbb{R} \to [0, 1]\) by

\begin{equation*} f_t = \chi_{\{t\}} \ \text{ if } t \in [0,1], \qquad f_t = 0 \ \text{ if } t \notin [0,1], \end{equation*}

where \(\chi_A\) denotes the characteristic function of \(A\) (so \(\chi_A(x) = 1\) if \(x \in A\) and \(\chi_A(x) = 0\) otherwise). Each \(f_t\) maps \(\mathbb{R}\) into \(\{0,1\} \subseteq [0,1]\), as required.

Each \(f_t\) is \(\mathcal{S}\)-measurable. If \(t \notin [0,1]\) then \(f_t\) is the constant function \(0\), and the inverse image under \(f_t\) of any set is either \(\emptyset\) or \(\mathbb{R}\), both of which lie in \(\mathcal{S}\). If \(t \in [0,1]\) then for a Borel set \(B \subseteq \mathbb{R}\) the inverse image \(f_t^{-1}(B)\) is one of the four sets \(\emptyset\), \(\{t\}\), \(\mathbb{R} \setminus \{t\}\), \(\mathbb{R}\), according to which of \(0, 1\) belong to \(B\). The set \(\{t\}\) is countable, and \(\mathbb{R} \setminus \{t\}\) has countable complement, namely \(\{t\}\); also \(\emptyset\) is countable and \(\mathbb{R}\) has empty (hence countable) complement. So all four of these sets are in \(\mathcal{S}\), and \(f_t\) is \(\mathcal{S}\)-measurable.

Now compute the pointwise supremum. Fix \(x \in \mathbb{R}\).

If \(x \in [0,1]\), then \(f_x = \chi_{\{x\}}\), so \(f_x(x) = 1\); since every \(f_t\) takes values in \([0,1]\) we get \(f(x) = \sup\{f_t(x) : t \in \mathbb{R}\} = 1\).

If \(x \notin [0,1]\), then for each \(t \in \mathbb{R}\) we have \(f_t(x) = 0\): either \(t \notin [0,1]\), in which case \(f_t\) is identically \(0\); or \(t \in [0,1]\), in which case \(t \neq x\) (because \(x \notin [0,1]\)) and so \(\chi_{\{t\}}(x) = 0\). Hence \(f(x) = 0\).

Therefore \(f = \chi_{[0,1]}\). This function is not \(\mathcal{S}\)-measurable, because

\begin{equation*} f^{-1}\big((\tfrac{1}{2}, \infty)\big) = [0,1], \end{equation*}

and \([0,1]\) is not in \(\mathcal{S}\): the interval \([0,1]\) is uncountable, and its complement \((-\infty, 0) \cup (1, \infty)\) is uncountable as well, since it contains the interval \((1,2)\). (Every open interval \((a,b)\) with \(a < b\) is uncountable: the map \(s \mapsto a + (b-a)\frac{e^{s}}{1 + e^{s}}\) is a bijection from \(\mathbb{R}\) onto \((a,b)\), and \(\mathbb{R}\) is uncountable by Cantor’s theorem. Hence \([0,1] \supseteq (0,1)\) and \((-\infty,0) \cup (1,\infty) \supseteq (1,2)\) are both uncountable.) Since \((\tfrac{1}{2}, \infty)\) is a Borel subset of \(\mathbb{R}\) whose inverse image under \(f\) is not in \(\mathcal{S}\), the function \(f\) fails to be \(\mathcal{S}\)-measurable.

Thus \(\{f_t\}_{t \in \mathbb{R}}\) is a family of \(\mathcal{S}\)-measurable functions from \(X\) to \([0,1]\) whose pointwise supremum is not \(\mathcal{S}\)-measurable.

Note where the argument in 2.53 breaks down. There one writes, for \(a \in \mathbb{R}\),

\begin{equation*} h^{-1}\big((a, \infty]\big) = \bigcup_{k=1}^{\infty} f_k^{-1}\big((a, \infty]\big), \end{equation*}

and this union is in \(\mathcal{S}\) because \(\sigma\)-algebras are closed under countable unions. With an index set of cardinality \(|\mathbb{R}|\) the corresponding identity

\begin{equation*} f^{-1}\big((a, \infty)\big) = \bigcup_{t \in \mathbb{R}} f_t^{-1}\big((a, \infty)\big) \end{equation*}

is still true, but it is an uncountable union, and \(\sigma\)-algebras need not be closed under such unions. The example above realizes exactly this failure: each \(f_t^{-1}\big((\tfrac{1}{2},\infty)\big)\) is a singleton or empty, hence in \(\mathcal{S}\), but the union of the singletons \(\{t\}\) over \(t \in [0,1]\) is the non-measurable set \([0,1]\).

Problem (2B.30)

Show that

\begin{equation*} \lim_{j \to \infty} \left( \lim_{k \to \infty} \big( \cos(j!\,\pi x) \big)^{2k} \right) = \begin{cases} 1 & \text{if } x \text{ is rational},\\ 0 & \text{if } x \text{ is irrational} \end{cases} \end{equation*}

for every \(x \in \mathbb{R}\).

[This example is due to Henri Lebesgue.]

Solution

Fix \(x \in \mathbb{R}\). The claim is proved in two steps: first the inner limit over \(k\) is evaluated for each fixed \(j \in \mathbb{Z}^{+}\), and then the outer limit over \(j\) is evaluated.

Step 1: the inner limit. Fix \(j \in \mathbb{Z}^{+}\) and set \(c = \cos(j!\,\pi x)\), so \(c \in [-1, 1]\) and \(c^{2k} = (c^2)^k\) with \(c^2 \in [0,1]\). Hence

\begin{equation*} \lim_{k \to \infty} c^{2k} = \lim_{k \to \infty} (c^2)^k = \begin{cases} 1 & \text{if } c^2 = 1,\\ 0 & \text{if } c^2 < 1, \end{cases} \end{equation*}

because \(1^k = 1\) for every \(k\), while \(0 \le r < 1\) implies \(r^k \to 0\). (The case \(c = 0\) is included: then \(c^{2k} = 0\) for every \(k \in \mathbb{Z}^{+}\).) Thus the inner limit exists for every \(j\) and every \(x\), and it equals \(1\) exactly when \(|\cos(j!\,\pi x)| = 1\).

Now \(|\cos \theta| = 1\) if and only if \(\theta\) is an integer multiple of \(\pi\). Applying this with \(\theta = j!\,\pi x\) gives

\begin{equation*} |\cos(j!\,\pi x)| = 1 \iff j!\,\pi x \in \{m\pi : m \in \mathbb{Z}\} \iff j!\,x \in \mathbb{Z}. \end{equation*}

Therefore, defining \(g_j : \mathbb{R} \to \mathbb{R}\) by \(g_j(x) = \lim_{k \to \infty} \big(\cos(j!\,\pi x)\big)^{2k}\), we have shown

\begin{equation*} g_j = \chi_{A_j}, \qquad \text{where } A_j = \{x \in \mathbb{R} : j!\,x \in \mathbb{Z}\} = \Big\{ \tfrac{m}{j!} : m \in \mathbb{Z} \Big\}, \end{equation*}

with \(\chi_{A_j}\) the characteristic function of \(A_j\).

Step 2: the outer limit. Suppose first that \(x\) is irrational. If we had \(j!\,x \in \mathbb{Z}\) for some \(j \in \mathbb{Z}^{+}\), then \(x = (j!\,x)/j!\) would be a quotient of integers with nonzero denominator, hence rational, a contradiction. Thus \(j!\,x \notin \mathbb{Z}\) for every \(j \in \mathbb{Z}^{+}\), so \(g_j(x) = 0\) for every \(j\), and consequently

\begin{equation*} \lim_{j \to \infty} g_j(x) = 0. \end{equation*}

Now suppose \(x\) is rational and write \(x = p/q\) with \(p \in \mathbb{Z}\) and \(q \in \mathbb{Z}^{+}\). If \(j \ge q\) then \(q\) is one of the factors in the product \(j! = 1 \cdot 2 \cdots q \cdots j\), so \(j!/q \in \mathbb{Z}\) and hence

\begin{equation*} j!\,x = \frac{j!\,p}{q} = p \cdot \frac{j!}{q} \in \mathbb{Z}. \end{equation*}

Thus \(g_j(x) = 1\) for every \(j \ge q\), and consequently

\begin{equation*} \lim_{j \to \infty} g_j(x) = 1. \end{equation*}

Combining the two cases,

\begin{equation*} \lim_{j \to \infty} \left( \lim_{k \to \infty} \big( \cos(j!\,\pi x) \big)^{2k} \right) = \begin{cases} 1 & \text{if } x \text{ is rational},\\ 0 & \text{if } x \text{ is irrational}, \end{cases} \end{equation*}

for every \(x \in \mathbb{R}\), as desired.

Remark. The double limit above is \(\chi_{\mathbb{Q}}\), whose restriction to \([0,1]\) is the function of Example 1.17 that is not Riemann integrable. Each function \(x \mapsto \big(\cos(j!\,\pi x)\big)^{2k}\) is continuous on \(\mathbb{R}\), hence Borel measurable by 2.41. Applying 2.48 twice, first in \(k\) and then in \(j\), shows again that \(\chi_{\mathbb{Q}}\) is Borel measurable. This illustrates the point made just before 2.48: unlike the Riemann integrable functions, the Borel measurable functions are closed under pointwise limits, so no iteration of such limits can escape measurability.

Exercises 2C

Problem (2C.1)

Explain why there does not exist a measure space \((X, \mathcal{S}, \mu)\) with the property that \(\{\mu(E) : E \in \mathcal{S}\} = [0, 1)\).

Solution

Suppose, seeking a contradiction, that \((X, \mathcal{S}, \mu)\) is a measure space with \[ \{\mu(E) : E \in \mathcal{S}\} = [0, 1). \]

Because \(\mathcal{S}\) is a \(\sigma\)-algebra on \(X\), we have \(X \in \mathcal{S}\). Hence \(\mu(X)\) belongs to the set above, and so \[ \mu(X) < 1. \]

Choose a real number \(t\) with \(\mu(X) < t < 1\); this is possible because \(\mu(X) < 1\). Then \(t \in [0, 1)\), so by our assumption there exists \(E \in \mathcal{S}\) with \(\mu(E) = t\).

However \(E \subseteq X\) and both sets are in \(\mathcal{S}\), so 2.57(a) gives \[ t = \mu(E) \le \mu(X) < t, \] which is a contradiction.

Thus no such measure space exists.

The essential point is that the set \(\{\mu(E) : E \in \mathcal{S}\}\) always contains the number \(\mu(X)\), which is an upper bound for that entire set. In other words, a bounded set of the form \(\{\mu(E) : E \in \mathcal{S}\}\) must contain its own supremum, and \([0,1)\) does not contain its supremum \(1\).

Problem (2C.2)

Let \(2^{\mathbb{Z}^+}\) denote the \(\sigma\)-algebra on \(\mathbb{Z}^+\) consisting of all subsets of \(\mathbb{Z}^+\).

Suppose \(\mu\) is a measure on \((\mathbb{Z}^+, 2^{\mathbb{Z}^+})\). Prove that there is a sequence \(w_1, w_2, \ldots\) in \([0, \infty]\) such that \[ \mu(E) = \sum_{k \in E} w_k \] for every set \(E \subseteq \mathbb{Z}^+\).

Solution

Define \[ w_k = \mu(\{k\}) \qquad \text{for } k \in \mathbb{Z}^+. \] Because \(\mu\) maps \(2^{\mathbb{Z}^+}\) into \([0, \infty]\), this is a sequence in \([0, \infty]\).

As in 2.55, for a subset \(E \subseteq \mathbb{Z}^+\) the sum \(\sum_{k \in E} w_k\) means the supremum of the finite subsums \(\sum_{k \in D} w_k\) as \(D\) ranges over the finite subsets of \(E\). We must show \(\mu(E)\) equals that supremum for every \(E \subseteq \mathbb{Z}^+\). Note that every subset of \(\mathbb{Z}^+\) lies in the \(\sigma\)-algebra \(2^{\mathbb{Z}^+}\), so \(\mu(E)\) is defined for every \(E\).

First suppose \(E = \emptyset\). Then \(\mu(E) = 0\) by the definition of a measure, and the only finite subset of \(E\) is \(\emptyset\), whose subsum is the empty sum \(0\). So both sides equal \(0\).

Next suppose \(E\) is a nonempty finite set, say \(E = \{k_1, \ldots, k_n\}\) with \(k_1 < \cdots < k_n\). The singletons \(\{k_1\}, \ldots, \{k_n\}\) are disjoint sets in \(2^{\mathbb{Z}^+}\) whose union is \(E\), so finite additivity (which follows from countable additivity, as noted after 2.54, by appending \(\emptyset, \emptyset, \ldots\) to the list) gives \[ \mu(E) = \sum_{j=1}^{n} \mu(\{k_j\}) = \sum_{j=1}^{n} w_{k_j}. \] Since \(E\) itself is the largest finite subset of \(E\) and all the \(w_k\) are nonnegative, the supremum defining \(\sum_{k \in E} w_k\) is attained at \(D = E\) and equals \(\sum_{j=1}^{n} w_{k_j}\). Hence \(\mu(E) = \sum_{k \in E} w_k\).

Finally suppose \(E\) is infinite. Because \(E \subseteq \mathbb{Z}^+\), we may list its elements in increasing order as \(k_1 < k_2 < \cdots\), and then \[ E = \bigcup_{j=1}^{\infty} \{k_j\} \] is a disjoint union of sets in \(2^{\mathbb{Z}^+}\). Countable additivity gives \[ \mu(E) = \sum_{j=1}^{\infty} w_{k_j} = \lim_{n \to \infty} \sum_{j=1}^{n} w_{k_j} = \sup_{n \in \mathbb{Z}^+} \sum_{j=1}^{n} w_{k_j}, \] where the last two equalities hold because the partial sums form an increasing sequence in \([0, \infty]\) (all terms are nonnegative), so the limit of the partial sums equals their supremum. (If some \(w_{k_j} = \infty\) or if the series diverges, both sides equal \(\infty\).)

It remains to see that \[ \sup_{n \in \mathbb{Z}^+} \sum_{j=1}^{n} w_{k_j} = \sup \Big\{ \sum_{k \in D} w_k : D \subseteq E \text{ and } D \text{ is finite} \Big\}. \] The left side is a supremum over the particular finite subsets \(D_n = \{k_1, \ldots, k_n\}\) of \(E\), so it is at most the right side. Conversely, if \(D \subseteq E\) is finite and nonempty, then \(D \subseteq D_n\) for \(n\) large enough (take \(n\) so large that \(D \subseteq \{k_1, \ldots, k_n\}\), which is possible since \(D\) is finite and \(k_1, k_2, \ldots\) lists all of \(E\)); as all \(w_k \ge 0\) this gives \(\sum_{k \in D} w_k \le \sum_{j=1}^{n} w_{k_j}\), which is at most the left side. Hence the right side is at most the left side, and the two suprema are equal.

Combining the displays, \(\mu(E) = \sum_{k \in E} w_k\) for every infinite \(E\) as well, completing the proof.

Problem (2C.3)

Give an example of a measure \(\mu\) on \((\mathbb{Z}^+, 2^{\mathbb{Z}^+})\) such that \[ \{\mu(E) : E \subseteq \mathbb{Z}^+\} = [0, 1]. \]

Solution

Define \(w : \mathbb{Z}^+ \to [0, \infty]\) by \(w(k) = 2^{-k}\) and let \[ \mu(E) = \sum_{k \in E} 2^{-k} \qquad \text{for } E \subseteq \mathbb{Z}^+. \] By the third bullet point of 2.55, \(\mu\) is a measure on \((\mathbb{Z}^+, 2^{\mathbb{Z}^+})\). We claim that \(\{\mu(E) : E \subseteq \mathbb{Z}^+\} = [0, 1]\).

One inclusion is easy. For every \(E \subseteq \mathbb{Z}^+\) we have \[ 0 \le \mu(E) \le \mu(\mathbb{Z}^+) = \sum_{k=1}^{\infty} 2^{-k} = 1, \] using 2.57(a) for the middle inequality. Hence \(\{\mu(E) : E \subseteq \mathbb{Z}^+\} \subseteq [0, 1]\).

For the reverse inclusion, let \(t \in [0, 1]\); we produce \(E \subseteq \mathbb{Z}^+\) with \(\mu(E) = t\). If \(t = 1\), take \(E = \mathbb{Z}^+\), which works by the computation above. So assume \(0 \le t < 1\).

Define \(b_1, b_2, \ldots\) in \(\{0, 1\}\) recursively as follows. Having defined \(b_1, \ldots, b_{k-1}\), set \(s_{k-1} = \sum_{j=1}^{k-1} b_j 2^{-j}\) (so \(s_0 = 0\)) and put \[ b_k = \begin{cases} 1 & \text{if } s_{k-1} + 2^{-k} \le t, \\ 0 & \text{otherwise.} \end{cases} \]

We prove by induction on \(k \ge 0\) that \[ s_k \le t < s_k + 2^{-k}. \] For \(k = 0\) this says \(0 \le t < 1\), which holds by assumption. Suppose the inequalities hold for \(k-1\), so \(s_{k-1} \le t < s_{k-1} + 2^{-(k-1)}\). There are two cases.

If \(s_{k-1} + 2^{-k} \le t\), then \(b_k = 1\) and \(s_k = s_{k-1} + 2^{-k}\), so \(s_k \le t\); moreover \[ t < s_{k-1} + 2^{-(k-1)} = s_{k-1} + 2^{-k} + 2^{-k} = s_k + 2^{-k}. \]

If \(s_{k-1} + 2^{-k} > t\), then \(b_k = 0\) and \(s_k = s_{k-1}\), so \(s_k = s_{k-1} \le t\) and \(t < s_{k-1} + 2^{-k} = s_k + 2^{-k}\).

In both cases the inequalities hold for \(k\), completing the induction.

From \(s_k \le t < s_k + 2^{-k}\) we get \(0 \le t - s_k < 2^{-k}\) for every \(k\), hence \(s_k \to t\) as \(k \to \infty\).

Now let \(E = \{k \in \mathbb{Z}^+ : b_k = 1\}\), so that \(s_k = \sum_{j \in E, \, j \le k} 2^{-j}\) for each \(k\); that is, \(s_k\) is the subsum of \(\sum_{k \in E} 2^{-k}\) over the finite set \(E \cap \{1, \ldots, k\}\). We check that \(\mu(E) = \sup_k s_k\). Each \(s_k\) is a subsum of \(\mu(E)\) over a finite subset of \(E\), so \(s_k \le \mu(E)\) and hence \(\sup_k s_k \le \mu(E)\). Conversely, if \(D \subseteq E\) is finite then \(D \subseteq E \cap \{1, \ldots, k\}\) for \(k\) large enough, so \(\sum_{j \in D} 2^{-j} \le s_k \le \sup_k s_k\) because all terms are nonnegative; taking the supremum over such \(D\) gives \(\mu(E) \le \sup_k s_k\). Since \(s_1 \le s_2 \le \cdots\), the supremum equals the limit, and therefore \[ \mu(E) = \sum_{k \in E} 2^{-k} = \lim_{k \to \infty} s_k = t. \]

Therefore every \(t \in [0,1]\) is of the form \(\mu(E)\), and \(\{\mu(E) : E \subseteq \mathbb{Z}^+\} = [0, 1]\).

Problem (2C.4)

Give an example of a measure space \((X, \mathcal{S}, \mu)\) such that \[ \{\mu(E) : E \in \mathcal{S}\} = \{\infty\} \cup \bigcup_{k=0}^{\infty} [3k, 3k+1]. \]

Solution

Let \(X = \mathbb{Z}^+ \times \{1, 2\}\), let \(\mathcal{S}\) be the \(\sigma\)-algebra of all subsets of \(X\), and define \(w : X \to [0, \infty]\) by \[ w(k, 1) = 2^{-k}, \qquad w(k, 2) = 3 \qquad \text{for } k \in \mathbb{Z}^+. \] Define \[ \mu(E) = \sum_{x \in E} w(x) \qquad \text{for } E \in \mathcal{S}, \] where as in 2.55 the sum is the supremum of the finite subsums. By the third bullet point of 2.55, \(\mu\) is a measure on \((X, \mathcal{S})\). Informally: \(X\) carries a part of total mass \(1\) whose measures fill out all of \([0,1]\) (the points \((k,1)\), with the weights of Exercise 3), together with infinitely many atoms of mass \(3\) (the points \((k,2)\)).

Given \(E \subseteq X\), write \[ E_1 = \{k \in \mathbb{Z}^+ : (k, 1) \in E\}, \qquad E_2 = \{k \in \mathbb{Z}^+ : (k, 2) \in E\}. \] Then \(E\) is the disjoint union of \(E_1 \times \{1\}\) and \(E_2 \times \{2\}\), both of which lie in \(\mathcal{S}\), so finite additivity gives \[ \mu(E) = s + b, \qquad \text{where } s = \sum_{k \in E_1} 2^{-k} \text{ and } b = \sum_{k \in E_2} 3. \] Here \(b = 3n\) if \(E_2\) is a finite set with exactly \(n\) elements, and \(b = \infty\) if \(E_2\) is infinite (in that case the finite subsums \(3, 6, 9, \ldots\) are unbounded, so their supremum is \(\infty\)).

Exercise 3 above showed that as \(E_1\) ranges over all subsets of \(\mathbb{Z}^+\), the number \(\sum_{k \in E_1} 2^{-k}\) ranges over exactly \([0, 1]\); so \(s\) ranges over exactly \([0,1]\). By the previous paragraph, as \(E_2\) ranges over all subsets of \(\mathbb{Z}^+\) the number \(b\) ranges over exactly \(\{0, 3, 6, 9, \ldots\} \cup \{\infty\}\). Moreover \(E_1\) and \(E_2\) can be chosen independently of each other, since every pair of subsets \((E_1, E_2)\) arises from the set \(E = (E_1 \times \{1\}) \cup (E_2 \times \{2\})\).

Hence \[ \{\mu(E) : E \in \mathcal{S}\} = \{s + b : s \in [0,1], \; b \in \{0, 3, 6, \ldots\} \cup \{\infty\}\}. \] If \(b = \infty\) the sum is \(\infty\) (for every \(s \in [0,1]\)). If \(b = 3k\) for some integer \(k \ge 0\), then as \(s\) runs over \([0,1]\) the sum \(3k + s\) runs over exactly \([3k, 3k+1]\). Therefore \[ \{\mu(E) : E \in \mathcal{S}\} = \{\infty\} \cup \bigcup_{k=0}^{\infty} [3k, 3k+1], \] as required.

Problem (2C.5)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space such that \(\mu(X) < \infty\). Prove that if \(\mathcal{A}\) is a set of disjoint sets in \(\mathcal{S}\) such that \(\mu(A) > 0\) for every \(A \in \mathcal{A}\), then \(\mathcal{A}\) is a countable set.

Solution

Here \(\mathcal{A}\) being a set of disjoint sets means that its elements are pairwise disjoint: if \(A, A^{\prime} \in \mathcal{A}\) and \(A \ne A^{\prime}\), then \(A \cap A^{\prime} = \emptyset\).

For each \(n \in \mathbb{Z}^+\) define \[ \mathcal{A}_n = \Big\{ A \in \mathcal{A} : \mu(A) > \tfrac{1}{n} \Big\}. \]

Step 1: each \(\mathcal{A}_n\) is a finite set, with at most \(n\,\mu(X)\) elements.

Suppose \(A_1, \ldots, A_m\) are \(m\) distinct elements of \(\mathcal{A}_n\). They are pairwise disjoint sets in \(\mathcal{S}\), so their union \(A_1 \cup \cdots \cup A_m\) lies in \(\mathcal{S}\) and finite additivity (which follows from countable additivity, as noted after 2.54) gives \[ \mu(A_1 \cup \cdots \cup A_m) = \sum_{j=1}^{m} \mu(A_j) > \frac{m}{n}. \] On the other hand \(A_1 \cup \cdots \cup A_m \subseteq X\), so 2.57(a) gives \[ \mu(A_1 \cup \cdots \cup A_m) \le \mu(X). \] Combining the two displays yields \(m/n < \mu(X)\), that is, \(m < n\,\mu(X)\). Since \(\mu(X) < \infty\), the number \(n\,\mu(X)\) is a finite real number. Thus every finite list of distinct elements of \(\mathcal{A}_n\) has length less than \(n\,\mu(X)\). If \(\mathcal{A}_n\) were an infinite set, we could choose \(m\) distinct elements of it for some integer \(m \ge n\,\mu(X)\), contradicting the previous sentence. Hence \(\mathcal{A}_n\) is finite.

Step 2: \(\mathcal{A} = \bigcup_{n=1}^{\infty} \mathcal{A}_n\).

Each \(\mathcal{A}_n \subseteq \mathcal{A}\), giving one inclusion. Conversely, if \(A \in \mathcal{A}\) then \(\mu(A) > 0\) by hypothesis, so by the Archimedean property there exists \(n \in \mathbb{Z}^+\) with \(\frac{1}{n} < \mu(A)\), and then \(A \in \mathcal{A}_n\). (Note \(\mu(A)\) could be \(\infty\), which also satisfies \(\mu(A) > \frac{1}{n}\) for every \(n\); in fact \(\mu(A) \le \mu(X) < \infty\) here, but we do not need this.)

Step 3: conclusion.

By Steps 1 and 2, \(\mathcal{A}\) is a countable union of finite sets, and a countable union of countable sets is countable. Therefore \(\mathcal{A}\) is a countable set.

Problem (2C.6)

Find all \(c \in [3, \infty)\) such that there exists a measure space \((X, \mathcal{S}, \mu)\) with \[ \{\mu(E) : E \in \mathcal{S}\} = [0, 1] \cup [3, c]. \]

Solution

The answer is that \(c = 4\) is the only such number.

Part 1: if such a measure space exists, then \(c = 4\).

Suppose \(c \in [3, \infty)\) and \((X, \mathcal{S}, \mu)\) is a measure space with \[ R := \{\mu(E) : E \in \mathcal{S}\} = [0, 1] \cup [3, c]. \]

First we compute \(\mu(X)\). Every element of \(R\) is at most \(c\), and \(c \in R\) (because \(c \in [3, c]\)). Since \(E \subseteq X\) for every \(E \in \mathcal{S}\), 2.57(a) gives \(\mu(E) \le \mu(X)\) for every \(E \in \mathcal{S}\); taking \(E\) with \(\mu(E) = c\) shows \(c \le \mu(X)\). Also \(\mu(X) \in R\), so \(\mu(X) \le c\). Hence \[ \mu(X) = c < \infty. \]

Next, \(R\) is symmetric about \(c/2\): if \(t \in R\), say \(t = \mu(E)\) with \(E \in \mathcal{S}\), then \(X \setminus E \in \mathcal{S}\), and \(\mu(E) = t \le c < \infty\) since every element of \(R\) is at most \(c\). So 2.57(b), applied to the pair \(E \subseteq X\) of sets in \(\mathcal{S}\) (in the notation of 2.57, with \(D = E\) and with \(X\) playing the role of the larger set), gives \[ \mu(X \setminus E) = \mu(X) - \mu(E) = c - t. \] Thus we have the reflection property \[ t \in R \implies c - t \in R. \]

Apply the reflection property with \(t = 1 \in R\): we get \(c - 1 \in R = [0,1] \cup [3,c]\). Since \(c \ge 3\) we have \(c - 1 \ge 2 > 1\), so \(c - 1 \notin [0,1]\); therefore \(c - 1 \in [3, c]\), which forces \(c - 1 \ge 3\), that is, \[ c \ge 4. \]

Now suppose, seeking a contradiction, that \(c \ge 5\). Then \(3 \le c - 2 \le c\), so \(c - 2 \in [3, c] \subseteq R\). Applying the reflection property with \(t = c - 2\) gives \(c - (c-2) = 2 \in R\). But \(2 \notin [0,1]\) and \(2 \notin [3, c]\), so \(2 \notin R\), a contradiction. Hence \[ 4 \le c < 5. \]

Finally apply the reflection property with \(t = 3 \in R\): we get \(c - 3 \in R\). From \(4 \le c < 5\) we have \(1 \le c - 3 < 2\). The only element of \(R = [0,1] \cup [3,c]\) lying in \([1, 2)\) is \(1\), so \(c - 3 = 1\) and therefore \[ c = 4. \]

Part 2: for \(c = 4\) such a measure space exists.

Let \(X = \{0, 1, 2, 3, \ldots\}\), let \(\mathcal{S}\) be the \(\sigma\)-algebra of all subsets of \(X\), and define \(w : X \to [0, \infty]\) by \[ w(0) = 3, \qquad w(k) = 2^{-k} \text{ for } k \in \mathbb{Z}^+. \] Let \(\mu(E) = \sum_{x \in E} w(x)\) for \(E \subseteq X\); by the third bullet point of 2.55 this is a measure on \((X, \mathcal{S})\).

For \(E \subseteq X\) write \(E^{\prime} = E \cap \mathbb{Z}^+\). Splitting \(E\) into the disjoint pieces \(E \cap \{0\}\) and \(E^{\prime}\) and using finite additivity, \[ \mu(E) = \begin{cases} \displaystyle\sum_{k \in E^{\prime}} 2^{-k} & \text{if } 0 \notin E, \\[2mm] 3 + \displaystyle\sum_{k \in E^{\prime}} 2^{-k} & \text{if } 0 \in E. \end{cases} \] By Exercise 3, as \(E^{\prime}\) ranges over all subsets of \(\mathbb{Z}^+\) the quantity \(\sum_{k \in E^{\prime}} 2^{-k}\) ranges over exactly \([0, 1]\), and the choice of \(E^{\prime}\) is independent of whether \(0 \in E\). Hence \[ \{\mu(E) : E \in \mathcal{S}\} = [0, 1] \cup \big(3 + [0,1]\big) = [0, 1] \cup [3, 4], \] which is \([0,1] \cup [3,c]\) with \(c = 4\).

Problem (2C.7)

Give an example of a measure space \((X, \mathcal{S}, \mu)\) such that \[ \{\mu(E) : E \in \mathcal{S}\} = [0, 1] \cup [3, \infty]. \]

Solution

Let \(X = \mathbb{Z}^+ \times \{1, 2\}\), let \(\mathcal{S}\) be the \(\sigma\)-algebra of all subsets of \(X\), and define \(w : X \to [0, \infty]\) by \[ w(k, 1) = 2^{-k}, \qquad w(k, 2) = k + 2 \qquad \text{for } k \in \mathbb{Z}^+. \] Define \[ \mu(E) = \sum_{x \in E} w(x) \qquad \text{for } E \in \mathcal{S}, \] the sum again meaning the supremum of the finite subsums, as in 2.55; by the third bullet point of 2.55, \(\mu\) is a measure on \((X, \mathcal{S})\).

The idea: the points \((k,1)\) carry the weights of Exercise 3, giving a part whose measures fill out exactly \([0,1]\); the points \((k,2)\) are atoms of masses \(3, 4, 5, 6, \ldots\), each of which is at least \(3\) (so no value in the forbidden gap \((1,3)\) can be produced) while the integers \(3, 4, 5, \ldots\) are spaced only \(1\) apart (so together with the \([0,1]\) part they sweep out all of \([3, \infty)\)).

Given \(E \subseteq X\), set \[ E_1 = \{k \in \mathbb{Z}^+ : (k,1) \in E\}, \qquad E_2 = \{k \in \mathbb{Z}^+ : (k,2) \in E\}. \] Since \(E\) is the disjoint union of \(E_1 \times \{1\}\) and \(E_2 \times \{2\}\), finite additivity gives \[ \mu(E) = s + b, \qquad \text{where } s = \sum_{k \in E_1} 2^{-k} \text{ and } b = \sum_{k \in E_2} (k+2). \] Note that \(E_1\) and \(E_2\) may be chosen independently: every pair of subsets \((E_1, E_2)\) of \(\mathbb{Z}^+\) arises from \(E = (E_1 \times \{1\}) \cup (E_2 \times \{2\})\).

By Exercise 3, as \(E_1\) ranges over all subsets of \(\mathbb{Z}^+\), the number \(s\) ranges over exactly \([0,1]\). Let \[ B = \Big\{ \sum_{k \in E_2} (k+2) : E_2 \subseteq \mathbb{Z}^+ \Big\} \] be the set of possible values of \(b\). We record three facts about \(B\).

(i) \(B \subseteq \{0\} \cup [3, \infty]\). Indeed, if \(E_2 = \emptyset\) then \(b = 0\). If \(E_2 \ne \emptyset\), pick \(k \in E_2\); then \(\{k\}\) is a finite subset of \(E_2\), so \(b \ge k + 2 \ge 3\).

(ii) Every integer \(n \ge 3\) belongs to \(B\): take \(E_2 = \{n - 2\}\), which is a subset of \(\mathbb{Z}^+\) since \(n - 2 \ge 1\), and then \(b = (n-2) + 2 = n\).

(iii) \(\infty \in B\): take \(E_2 = \mathbb{Z}^+\); the finite subsums \(\sum_{k=1}^{m}(k+2)\) are unbounded, so their supremum is \(\infty\).

Now we show \(\{\mu(E) : E \in \mathcal{S}\} = [0,1] \cup [3, \infty]\).

Containment \(\subseteq\): let \(E \in \mathcal{S}\) and write \(\mu(E) = s + b\) as above, with \(s \in [0,1]\) and \(b \in B\). By (i) either \(b = 0\), in which case \(\mu(E) = s \in [0,1]\); or \(b \ge 3\) (possibly \(b = \infty\)), in which case \(\mu(E) = s + b \ge 3\), so \(\mu(E) \in [3, \infty]\). Either way \(\mu(E) \in [0,1] \cup [3,\infty]\).

Containment \(\supseteq\): let \(t \in [0,1] \cup [3, \infty]\).

If \(t \in [0,1]\), choose \(E_1 \subseteq \mathbb{Z}^+\) with \(\sum_{k \in E_1} 2^{-k} = t\) (possible by Exercise 3) and take \(E_2 = \emptyset\); then \(\mu(E) = t\).

If \(t = \infty\), take \(E_1 = \emptyset\) and \(E_2 = \mathbb{Z}^+\); then \(\mu(E) = \infty\) by (iii).

If \(t \in [3, \infty)\), let \(n = \lfloor t \rfloor\), the greatest integer with \(n \le t\). Then \(n \ge 3\) and \(t - n \in [0, 1)\). Take \(E_2 = \{n-2\}\), so that \(b = n\) by (ii), and choose \(E_1 \subseteq \mathbb{Z}^+\) with \(s = \sum_{k \in E_1} 2^{-k} = t - n\) (possible by Exercise 3, since \(t - n \in [0,1]\)). Then for \(E = (E_1 \times \{1\}) \cup (E_2 \times \{2\})\) we get \[ \mu(E) = (t - n) + n = t. \]

Hence every element of \([0,1] \cup [3, \infty]\) is a value of \(\mu\), and combining the two containments, \[ \{\mu(E) : E \in \mathcal{S}\} = [0, 1] \cup [3, \infty], \] as required.

Problem (2C.8)

Give an example of a set \(X\), a \(\sigma\)-algebra \(\mathcal{S}\) of subsets of \(X\), a set \(\mathcal{A}\) of subsets of \(X\) such that the smallest \(\sigma\)-algebra on \(X\) containing \(\mathcal{A}\) is \(\mathcal{S}\), and two measures \(\mu\) and \(\nu\) on \((X, \mathcal{S})\) such that \(\mu(A) = \nu(A)\) for all \(A \in \mathcal{A}\) and \(\mu(X) = \nu(X) < \infty\), but \(\mu \neq \nu\).

Solution

Take

\begin{equation*} X = \{1, 2, 3, 4\}, \qquad \mathcal{S} = \{E : E \subseteq X\}, \qquad \mathcal{A} = \big\{\{1,2\}, \{2,3\}\big\}, \end{equation*}

and let \(\mu\) and \(\nu\) be the measures on \((X, \mathcal{S})\) defined by

\begin{equation*} \mu = \delta_2 + \delta_4, \qquad \nu = \delta_1 + \delta_3, \end{equation*}

where \(\delta_c\) denotes the Dirac measure at \(c\) (see 2.55). Each of \(\mu\) and \(\nu\) is a measure, being a sum of two measures (Exercise 9 of this section, or the direct verification that a weighted counting measure \(E \mapsto \sum_{x \in E} w(x)\) is a measure, as in 2.55). Explicitly, for \(E \subseteq X\),

\begin{equation*} \mu(E) = \#\big(E \cap \{2,4\}\big), \qquad \nu(E) = \#\big(E \cap \{1,3\}\big). \end{equation*}

First, \(\mathcal{S}\) really is the smallest \(\sigma\)-algebra on \(X\) containing \(\mathcal{A}\). Indeed, let \(\mathcal{T}\) be any \(\sigma\)-algebra on \(X\) with \(\mathcal{A} \subseteq \mathcal{T}\). Then \(\mathcal{T}\) is closed under finite intersections and relative complements (because it is closed under complementation and countable unions), so

\begin{equation*} \{2\} = \{1,2\} \cap \{2,3\} \in \mathcal{T}, \quad \{1\} = \{1,2\} \setminus \{2\} \in \mathcal{T}, \quad \{3\} = \{2,3\} \setminus \{2\} \in \mathcal{T}, \end{equation*}

and hence \(\{4\} = X \setminus (\{1\} \cup \{2\} \cup \{3\}) \in \mathcal{T}\). Thus \(\mathcal{T}\) contains every singleton of \(X\); as \(X\) is finite, every subset of \(X\) is a finite union of singletons, so \(\mathcal{T} = \mathcal{S}\). In particular the smallest \(\sigma\)-algebra containing \(\mathcal{A}\) is \(\mathcal{S}\).

Next, \(\mu\) and \(\nu\) agree on \(\mathcal{A}\):

\begin{equation*} \mu(\{1,2\}) = 1 = \nu(\{1,2\}), \qquad \mu(\{2,3\}) = 1 = \nu(\{2,3\}). \end{equation*}

They also agree on \(X\), with a finite common value:

\begin{equation*} \mu(X) = 2 = \nu(X) < \infty. \end{equation*}

Finally \(\mu \neq \nu\), since

\begin{equation*} \mu(\{1\}) = 0 \neq 1 = \nu(\{1\}). \end{equation*}

This example shows that agreement of two finite measures on a generating family \(\mathcal{A}\) (together with agreement of the total masses) does not force the measures to be equal. The standard uniqueness theorems require in addition that \(\mathcal{A}\) be closed under finite intersections; here that hypothesis fails, because \(\{1,2\} \cap \{2,3\} = \{2\} \notin \mathcal{A}\).

Problem (2C.9)

Suppose \(\mu\) and \(\nu\) are measures on a measurable space \((X, \mathcal{S})\). Prove that \(\mu + \nu\) is a measure on \((X, \mathcal{S})\). [Here \(\mu + \nu\) is the usual sum of two functions: if \(E \in \mathcal{S}\), then \((\mu + \nu)(E) = \mu(E) + \nu(E)\).]

Solution

Throughout we use the usual arithmetic of \([0, \infty]\): \(a + \infty = \infty + a = \infty\) for every \(a \in [0, \infty]\). With this convention \([0, \infty]\) is closed under addition, and addition on \([0,\infty]\) is commutative and associative.

Since \(\mu(E) \in [0,\infty]\) and \(\nu(E) \in [0,\infty]\) for each \(E \in \mathcal{S}\), the function \(\mu + \nu\) maps \(\mathcal{S}\) into \([0, \infty]\). Also

\begin{equation*} (\mu + \nu)(\varnothing) = \mu(\varnothing) + \nu(\varnothing) = 0 + 0 = 0 . \end{equation*}

It remains to prove countable additivity. Before doing so we record an elementary fact about series in \([0, \infty]\).

Claim. If \((a_k)\) and \((b_k)\) are sequences in \([0, \infty]\), then

\begin{equation*} \sum_{k=1}^{\infty} (a_k + b_k) = \sum_{k=1}^{\infty} a_k + \sum_{k=1}^{\infty} b_k . \end{equation*}

Here each such series of nonnegative terms is by definition the limit in \([0,\infty]\) of its increasing sequence of partial sums, equivalently the supremum of the partial sums (which exists in \([0,\infty]\) whether or not it is finite).

Proof of Claim. For each \(n \in \mathbf{Z}^{+}\) let

\begin{equation*} A_n = \sum_{k=1}^{n} a_k, \qquad B_n = \sum_{k=1}^{n} b_k, \qquad C_n = \sum_{k=1}^{n}(a_k + b_k). \end{equation*}

Rearranging a finite sum of elements of \([0,\infty]\) gives \(C_n = A_n + B_n\). The sequences \((A_n)\), \((B_n)\), \((C_n)\) are increasing in \([0, \infty]\), with suprema \(A\), \(B\), \(C\) say, and these suprema are the three series in question.

If \(A = \infty\), then for every \(M \in (0,\infty)\) there is \(n\) with \(A_n > M\), hence \(C_n \geq A_n > M\); thus \(C = \infty = A + B\). The case \(B = \infty\) is symmetric. So assume \(A < \infty\) and \(B < \infty\). Then \(C_n = A_n + B_n \leq A + B\) for all \(n\), so \(C \leq A + B\). Conversely, given \(\varepsilon > 0\) choose \(n\) with \(A_n > A - \varepsilon/2\) and \(B_n > B - \varepsilon/2\) (possible, using the same \(n\) for both by monotonicity); then \(C_n = A_n + B_n > A + B - \varepsilon\), so \(C > A + B - \varepsilon\). As \(\varepsilon > 0\) was arbitrary, \(C \geq A + B\). Hence \(C = A + B\), proving the claim.

Now let \(E_1, E_2, \dots\) be a disjoint sequence of sets in \(\mathcal{S}\). Using the countable additivity of \(\mu\) and of \(\nu\) and then the claim with \(a_k = \mu(E_k)\) and \(b_k = \nu(E_k)\), we get

\begin{equation*} \begin{aligned} (\mu + \nu)\Big(\bigcup_{k=1}^{\infty} E_k\Big) &= \mu\Big(\bigcup_{k=1}^{\infty} E_k\Big) + \nu\Big(\bigcup_{k=1}^{\infty} E_k\Big) \\ &= \sum_{k=1}^{\infty} \mu(E_k) + \sum_{k=1}^{\infty} \nu(E_k) \\ &= \sum_{k=1}^{\infty} \big(\mu(E_k) + \nu(E_k)\big) \\ &= \sum_{k=1}^{\infty} (\mu + \nu)(E_k). \end{aligned} \end{equation*}

Thus \(\mu + \nu : \mathcal{S} \to [0, \infty]\) satisfies \((\mu+\nu)(\varnothing) = 0\) and is countably additive, so \(\mu + \nu\) is a measure on \((X, \mathcal{S})\), as required.

The same argument shows more generally that \(c\mu + d\nu\) is a measure for all \(c, d \in [0, \infty)\), and (by an extra application of the interchange of two suprema) that any countable sum \(\sum_{j=1}^{\infty} \mu_j\) of measures on \((X, \mathcal{S})\) is again a measure on \((X, \mathcal{S})\).

Problem (2C.10)

Give an example of a measure space \((X, \mathcal{S}, \mu)\) and a decreasing sequence \(E_1 \supseteq E_2 \supseteq \cdots\) of sets in \(\mathcal{S}\) such that

\begin{equation*} \mu\Big(\bigcap_{k=1}^{\infty} E_k\Big) \neq \lim_{k \to \infty} \mu(E_k). \end{equation*}

Solution

Take \(X = \mathbf{Z}^{+}\), let \(\mathcal{S}\) be the \(\sigma\)-algebra of all subsets of \(\mathbf{Z}^{+}\), and let \(\mu\) be counting measure on \((X, \mathcal{S})\): thus \(\mu(E) = n\) if \(E\) is a finite set with exactly \(n\) elements, and \(\mu(E) = \infty\) if \(E\) is infinite. As noted in 2.55, counting measure is indeed a measure. Define

\begin{equation*} E_k = \{k, k+1, k+2, \dots\} \quad \text{for } k \in \mathbf{Z}^{+}. \end{equation*}

Then \(E_1 \supseteq E_2 \supseteq \cdots\) is a decreasing sequence of sets in \(\mathcal{S}\).

Each \(E_k\) is an infinite subset of \(\mathbf{Z}^{+}\), so \(\mu(E_k) = \infty\) for every \(k\), and therefore

\begin{equation*} \lim_{k \to \infty} \mu(E_k) = \infty . \end{equation*}

On the other hand

\begin{equation*} \bigcap_{k=1}^{\infty} E_k = \varnothing, \end{equation*}

because if \(n \in \mathbf{Z}^{+}\) then \(n \notin E_{n+1}\), so no positive integer lies in every \(E_k\). Hence

\begin{equation*} \mu\Big(\bigcap_{k=1}^{\infty} E_k\Big) = \mu(\varnothing) = 0 \neq \infty = \lim_{k \to \infty} \mu(E_k), \end{equation*}

as desired.

This example shows that the hypothesis \(\mu(E_1) < \infty\) in 2.60 cannot be deleted. In fact that hypothesis fails as badly as possible here, and necessarily so: if \((E_k)\) is a decreasing sequence in \(\mathcal{S}\) with \(\mu(E_n) < \infty\) for some \(n \in \mathbf{Z}^{+}\), then applying 2.60 to the decreasing sequence \(E_n \supseteq E_{n+1} \supseteq \cdots\) (whose intersection equals \(\bigcap_{k=1}^{\infty} E_k\)) gives

\begin{equation*} \mu\Big(\bigcap_{k=1}^{\infty} E_k\Big) = \lim_{k \to \infty} \mu(E_k), \end{equation*}

where we used that discarding the first \(n - 1\) terms does not change the limit. Thus any counterexample must have \(\mu(E_k) = \infty\) for every \(k\).

Problem (2C.11)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(C, D, E \in \mathcal{S}\) are such that

\begin{equation*} \mu(C \cap D) < \infty, \qquad \mu(C \cap E) < \infty, \qquad \mu(D \cap E) < \infty . \end{equation*}

Find and prove a formula for \(\mu(C \cup D \cup E)\) in terms of \(\mu( C)\), \(\mu(D)\), \(\mu(E)\), \(\mu(C \cap D)\), \(\mu(C \cap E)\), \(\mu(D \cap E)\), and \(\mu(C \cap D \cap E)\).

Solution

The formula is the inclusion-exclusion identity for three sets:

\begin{equation*} \mu(C \cup D \cup E) = \mu( C) + \mu(D) + \mu(E) - \mu(C \cap D) - \mu(C \cap E) - \mu(D \cap E) + \mu(C \cap D \cap E). \end{equation*}

Notation and well-definedness. Write

\begin{equation*} p = \mu(C \cap D), \quad q = \mu(C \cap E), \quad r = \mu(D \cap E), \quad s = \mu(C \cap D \cap E). \end{equation*}

By hypothesis \(p, q, r < \infty\), and since \(C \cap D \cap E \subseteq C \cap D\) we get \(s \leq p < \infty\) from 2.57(a). So all four of \(p, q, r, s\) are (finite) nonnegative real numbers, while \(\mu( C), \mu(D), \mu(E)\) may be \(\infty\). The right side of the displayed formula is therefore unambiguous: it is an element of \([0,\infty]\) from which only finite quantities are subtracted, with the convention \(\infty - t = \infty\) for \(t \in \mathbf{R}\). To avoid any subtraction at all in the proof, we shall establish the equivalent identity

\begin{equation*} \mu(C \cup D \cup E) + p + q + r = \mu( C) + \mu(D) + \mu(E) + s , \end{equation*}

in which every term lies in \([0, \infty]\) and only addition occurs; since \(p, q, r, s\) are finite reals, this identity is equivalent to the displayed formula above.

Step 1. Since \(\mu(C \cap D) = p < \infty\), the result 2.61 applied to \(C\) and \(D\) gives \(\mu(C \cup D) = \mu( C) + \mu(D) - p\), that is,

\begin{equation*} \mu(C \cup D) + p = \mu( C) + \mu(D). \end{equation*}

Step 2. Note that

\begin{equation*} (C \cup D) \cap E = (C \cap E) \cup (D \cap E), \end{equation*}

and that \((C \cap E) \cap (D \cap E) = C \cap D \cap E\), which has measure \(s < \infty\). Hence 2.61 applies to the sets \(C \cap E\) and \(D \cap E\) and yields

\begin{equation*} \mu\big((C \cup D) \cap E\big) = q + r - s . \end{equation*}

In particular \(\mu\big((C \cup D) \cap E\big) \leq q + r < \infty\), and

\begin{equation*} \mu\big((C \cup D) \cap E\big) + s = q + r . \end{equation*}

Step 3. Because \(\mu\big((C \cup D) \cap E\big) < \infty\) by Step 2, the result 2.61 applies to the sets \(C \cup D\) and \(E\) and gives

\begin{equation*} \mu(C \cup D \cup E) = \mu(C \cup D) + \mu(E) - \mu\big((C \cup D) \cap E\big), \end{equation*}

which, since the subtracted quantity is finite, is the same as

\begin{equation*} \mu(C \cup D \cup E) + \mu\big((C \cup D) \cap E\big) = \mu(C \cup D) + \mu(E). \end{equation*}

Step 4. Add \(s\) to both sides of the identity in Step 3 and use Step 2 to replace \(\mu\big((C \cup D) \cap E\big) + s\) by \(q + r\):

\begin{equation*} \mu(C \cup D \cup E) + q + r = \mu(C \cup D) + \mu(E) + s . \end{equation*}

Now add \(p\) to both sides and use Step 1 to replace \(\mu(C \cup D) + p\) by \(\mu( C) + \mu(D)\):

\begin{equation*} \mu(C \cup D \cup E) + p + q + r = \mu( C) + \mu(D) + \mu(E) + s . \end{equation*}

(All rearrangements are legitimate: addition on \([0, \infty]\) is commutative and associative.)

This is exactly the addition-only identity announced above. Finally, since \(p\), \(q\), \(r\), \(s\) are finite real numbers, we may subtract \(p + q + r\) from both sides to conclude

\begin{equation*} \mu(C \cup D \cup E) = \mu( C) + \mu(D) + \mu(E) - \mu(C \cap D) - \mu(C \cap E) - \mu(D \cap E) + \mu(C \cap D \cap E), \end{equation*}

as claimed. (If \(\mu( C)\), \(\mu(D)\), or \(\mu(E)\) equals \(\infty\), both sides equal \(\infty\): the left side because \(\mu(C \cup D \cup E) \geq \mu( C)\) and similarly for \(D\) and \(E\) by 2.57(a), and the right side because a finite quantity is being subtracted from \(\infty\).)

Problem (2C.12)

Suppose \(X\) is a set and \(\mathcal{S}\) is the \(\sigma\)-algebra of all subsets \(E\) of \(X\) such that \(E\) is countable or \(X \setminus E\) is countable. Give a complete description of the set of all measures on \((X, \mathcal{S})\).

Solution

Notation. For a function \(w : X \to [0, \infty]\) and a set \(A \subseteq X\) put

\begin{equation*} \sum_{x \in A} w(x) = \sup\Big\{ \sum_{x \in D} w(x) : D \subseteq A \text{ and } D \text{ is finite} \Big\} \in [0, \infty], \end{equation*}

as in the third bullet of 2.55 (the empty sum is \(0\)). For \(E \in \mathcal{S}\) write \([E \text{ uncountable}]\) for \(1\) if \(E\) is uncountable and \(0\) if \(E\) is countable. Products of elements of \([0,\infty]\) are computed with the usual convention \(0 \cdot \infty = 0\); thus \(c \cdot [E \text{ uncountable}] = 0\) whenever \(E\) is countable, even if \(c = \infty\).

Answer. The measures on \((X, \mathcal{S})\) are exactly the functions \(\mu_{w,c} : \mathcal{S} \to [0,\infty]\) of the form

\begin{equation*} \mu_{w,c}(E) = \sum_{x \in E} w(x) \; + \; c \cdot [E \text{ uncountable}], \end{equation*}

where \(w : X \to [0, \infty]\) is an arbitrary function and \(c \in [0, \infty]\) is an arbitrary constant. In other words: a measure here is a weighted counting measure with weights \(w(x) = \mu(\{x\})\), plus an extra lump of mass \(c\) carried by every uncountable set of \(\mathcal{S}\) (when \(X\) is uncountable, the uncountable sets in \(\mathcal{S}\) are exactly the sets in \(\mathcal{S}\) with countable complement). If \(X\) itself is countable, then \(\mathcal{S}\) is the \(\sigma\)-algebra of all subsets of \(X\), no set is uncountable, and the description collapses to: the measures are exactly \(E \mapsto \sum_{x \in E} w(x)\) with \(w : X \to [0,\infty]\) arbitrary.

We prove three lemmas, then verify each \(\mu_{w,c}\) is a measure, then show every measure is of this form.

Lemma 1 (unordered sums are countably additive). If \(w : X \to [0,\infty]\) and \(A = \bigcup_{k=1}^{\infty} A_k\) with \(A_1, A_2, \dots\) disjoint subsets of \(X\), then \(\sum_{x \in A} w(x) = \sum_{k=1}^{\infty} \sum_{x \in A_k} w(x)\).

Proof. Let \(D \subseteq A\) be finite. Then \(D\) meets only finitely many \(A_k\), say only those with \(k \leq n\), and \(D\) is the disjoint union of the sets \(D \cap A_k\) for \(k \leq n\); hence

\begin{equation*} \sum_{x \in D} w(x) = \sum_{k=1}^{n} \sum_{x \in D \cap A_k} w(x) \leq \sum_{k=1}^{n} \sum_{x \in A_k} w(x) \leq \sum_{k=1}^{\infty} \sum_{x \in A_k} w(x). \end{equation*}

Taking the supremum over such \(D\) gives \(\sum_{x \in A} w(x) \leq \sum_{k=1}^{\infty} \sum_{x \in A_k} w(x)\). Conversely, fix \(n \in \mathbf{Z}^{+}\) and finite sets \(D_k \subseteq A_k\) for \(k = 1, \dots, n\). Then \(D_1 \cup \cdots \cup D_n\) is a finite subset of \(A\), so

\begin{equation*} \sum_{k=1}^{n} \sum_{x \in D_k} w(x) = \sum_{x \in D_1 \cup \cdots \cup D_n} w(x) \leq \sum_{x \in A} w(x). \end{equation*}

Taking suprema over \(D_1\), then \(D_2\), and so on (finitely many independent suprema, each legitimate because addition on \([0,\infty]\) is order preserving and commutes with suprema of increasing families in one variable) gives \(\sum_{k=1}^{n} \sum_{x \in A_k} w(x) \leq \sum_{x \in A} w(x)\). Letting \(n \to \infty\) finishes the proof.

Lemma 2. If \(N \subseteq X\) is uncountable and \(w(x) > 0\) for every \(x \in N\), then \(\sum_{x \in N} w(x) = \infty\).

Proof. We have \(N = \bigcup_{n=1}^{\infty} \{x \in N : w(x) > 1/n\}\). A countable union of countable sets is countable, so some \(\{x \in N : w(x) > 1/n\}\) is uncountable, in particular infinite. For each \(m \in \mathbf{Z}^{+}\) that set therefore contains a finite subset \(D\) with \(m\) elements, and \(\sum_{x \in D} w(x) > m/n\). Hence the supremum defining \(\sum_{x \in N} w(x)\) is \(\infty\).

Lemma 3. Let \(E_1, E_2, \dots\) be a disjoint sequence of sets in \(\mathcal{S}\) and \(E = \bigcup_{k=1}^{\infty} E_k\). Then \(E\) is uncountable if and only if exactly one \(E_k\) is uncountable.

Proof. First, at most one \(E_k\) can be uncountable: if \(E_i\) and \(E_j\) were uncountable with \(i \neq j\), then (being in \(\mathcal{S}\) and uncountable) each has countable complement, and \(E_i \subseteq X \setminus E_j\) because \(E_i \cap E_j = \varnothing\); thus \(E_i\) is a subset of a countable set, contradicting that \(E_i\) is uncountable. Next, if every \(E_k\) is countable then \(E\) is countable, being a countable union of countable sets. Conversely, if some \(E_k\) is uncountable then \(E \supseteq E_k\) is uncountable. Combining these statements gives the lemma.

Step 1: every \(\mu_{w,c}\) is a measure on \((X, \mathcal{S})\).

Clearly \(\mu_{w,c}\) maps \(\mathcal{S}\) into \([0, \infty]\), and \(\mu_{w,c}(\varnothing) = 0 + 0 = 0\) because \(\varnothing\) is countable and the empty sum is \(0\). Let \(E_1, E_2, \dots\) be a disjoint sequence in \(\mathcal{S}\) and \(E = \bigcup_{k=1}^{\infty} E_k\). By Lemma 3,

\begin{equation*} [E \text{ uncountable}] = \sum_{k=1}^{\infty} [E_k \text{ uncountable}], \end{equation*}

since either every term on the right is \(0\) and so is the left side, or exactly one term on the right is \(1\) and the left side is \(1\). Because at most one term of that series is nonzero, we may multiply through by \(c\) termwise with no ambiguity even when \(c = \infty\):

\begin{equation*} c\,[E \text{ uncountable}] = \sum_{k=1}^{\infty} c\,[E_k \text{ uncountable}], \end{equation*}

both sides being \(c\) if some (necessarily unique) \(E_k\) is uncountable and \(0\) otherwise. Combining this with Lemma 1 for the weight part, we get

\begin{equation*} \mu_{w,c}(E) = \sum_{x \in E} w(x) + c\,[E \text{ uncountable}] = \sum_{k=1}^{\infty} \Big( \sum_{x \in E_k} w(x) + c\,[E_k \text{ uncountable}] \Big) = \sum_{k=1}^{\infty} \mu_{w,c}(E_k), \end{equation*}

where the middle equality also uses that two series with terms in \([0,\infty]\) may be added termwise (proved in the solution to Exercise 9 of this section). Hence \(\mu_{w,c}\) is countably additive, so it is a measure.

Step 2: every measure on \((X, \mathcal{S})\) has this form.

Let \(\mu\) be a measure on \((X, \mathcal{S})\) and define

\begin{equation*} w(x) = \mu(\{x\}) \quad \text{for } x \in X, \qquad N = \{x \in X : w(x) > 0\}. \end{equation*}

(Each singleton is countable, hence in \(\mathcal{S}\).)

First, \(\mu(E) = \sum_{x \in E} w(x)\) for every countable \(E \in \mathcal{S}\). Indeed, list the distinct elements of \(E\) as \(x_1, x_2, \dots\) (a finite or infinite list) and apply countable additivity to the disjoint sequence \(\{x_1\}, \{x_2\}, \dots\) (padded with copies of \(\varnothing\) if \(E\) is finite) to get \(\mu(E) = \sum_k w(x_k)\); for nonnegative terms this series equals the unordered sum \(\sum_{x \in E} w(x)\), since its partial sums are exactly the subsums over the finite sets \(\{x_1, \dots, x_n\}\) and these are cofinal among all finite subsets of \(E\).

Second, note that for every \(E \in \mathcal{S}\) and every finite \(D \subseteq E\) we have \(\mu(E) \geq \mu(D) = \sum_{x \in D} w(x)\), using 2.57(a) and finite additivity; taking the supremum over \(D\) gives

\begin{equation*} \mu(E) \geq \sum_{x \in E} w(x) \qquad \text{for all } E \in \mathcal{S}. \end{equation*}

Now distinguish two cases.

Case 1: \(N\) is uncountable. Let \(G \in \mathcal{S}\) be uncountable; then \(X \setminus G\) is countable, so \(G \cap N = N \setminus (X \setminus G)\) is uncountable. By Lemma 2, \(\sum_{x \in G} w(x) \geq \sum_{x \in G \cap N} w(x) = \infty\), and by the inequality just proved \(\mu(G) = \infty\) as well. Hence \(\mu(G) = \sum_{x \in G} w(x) + 0\) for uncountable \(G\), while \(\mu(E) = \sum_{x \in E} w(x) + 0\) for countable \(E\) by the first paragraph of this step. Thus \(\mu = \mu_{w,0}\).

Case 2: \(N\) is countable. Then \(X \setminus N \in \mathcal{S}\); set

\begin{equation*} c = \mu(X \setminus N) \in [0, \infty]. \end{equation*}

Let \(G \in \mathcal{S}\) be uncountable, so \(X \setminus G\) is countable. The sets \(G \cap N\) and \(G \setminus N\) belong to \(\mathcal{S}\) (the first is countable; the second has complement \((X \setminus G) \cup N\), which is countable), and they are disjoint with union \(G\), so

\begin{equation*} \mu(G) = \mu(G \cap N) + \mu(G \setminus N). \end{equation*}

Since \(G \cap N\) is countable, \(\mu(G \cap N) = \sum_{x \in G \cap N} w(x) = \sum_{x \in G} w(x)\), the last equality because \(w\) vanishes off \(N\). For the other term, decompose \(X \setminus N\) as the disjoint union of \(G \setminus N\) and \((X \setminus N) \cap (X \setminus G)\); the latter set is countable (being a subset of \(X \setminus G\)) and \(w\) vanishes on it (being a subset of \(X \setminus N\)), so its measure is \(\sum_{x \in (X \setminus N) \cap (X \setminus G)} w(x) = 0\). Hence

\begin{equation*} c = \mu(X \setminus N) = \mu(G \setminus N) + 0 = \mu(G \setminus N). \end{equation*}

Combining the last three displays gives \(\mu(G) = \sum_{x \in G} w(x) + c\) for every uncountable \(G \in \mathcal{S}\), while \(\mu(E) = \sum_{x \in E} w(x)\) for countable \(E\). Thus \(\mu = \mu_{w,c}\).

In both cases \(\mu\) has the stated form, which completes the description.

Remarks on uniqueness. In the representation \(\mu = \mu_{w,c}\) the weight function is always uniquely determined, since each singleton is countable and hence \(w(x) = \mu_{w,c}(\{x\})\). The constant \(c\), however, need not be. Precisely, \(c\) is uniquely determined by \(\mu\) if and only if \(X\) is uncountable and \(N = \{x \in X : w(x) > 0\}\) is countable. Indeed, if \(X\) is countable then \(\mathcal{S}\) contains no uncountable set, so \(\mu_{w,c} = \mu_{w,c^{\prime}}\) for all \(c, c^{\prime} \in [0,\infty]\); and if \(N\) is uncountable then Lemma 2 (applied as in Case 1) gives \(\sum_{x \in G} w(x) = \infty\) for every uncountable \(G \in \mathcal{S}\), so again \(\mu_{w,c}\) does not depend on \(c\). In both of these degenerate situations one may normalize \(c = 0\). If instead \(X\) is uncountable and \(N\) is countable, then \(X \setminus N\) is an uncountable set in \(\mathcal{S}\) with \(\sum_{x \in X \setminus N} w(x) = 0\), so \(c = \mu(X \setminus N)\) is recovered from \(\mu\).

For example, taking \(w = 0\) and \(c = 3\) recovers the measure of the fourth bullet of 2.55, which assigns \(0\) to countable sets and \(3\) to uncountable sets in \(\mathcal{S}\).

Exercises 2D

Problem (2D.1)

(a) Show that the set consisting of those numbers in \((0,1)\) that have a decimal expansion containing one hundred consecutive 4s is a Borel subset of \(\mathbb{R}\).

(b) What is the Lebesgue measure of the set in (a)?

Solution

Throughout, a decimal expansion of \(x \in [0,1]\) means a sequence \(d_1, d_2, \ldots\) of digits in \(\{0,1,\ldots,9\}\) with \(x = \sum_{j=1}^{\infty} d_j 10^{-j}\). Let

\begin{equation*} E = \{x \in (0,1) : \text{some decimal expansion of } x \text{ contains one hundred consecutive } 4\text{s}\}. \end{equation*}

A basic observation. For a finite string \(s = (d_1,\ldots,d_m) \in \{0,\ldots,9\}^m\) put

\begin{equation*} a_s = \sum_{j=1}^{m} d_j 10^{-j}, \qquad I_s = [a_s,\, a_s + 10^{-m}]. \end{equation*}

Then a number \(x \in [0,1]\) has a decimal expansion beginning with the string \(s\) if and only if \(x \in I_s\).

Indeed, if \(x = \sum_{j=1}^\infty d_j 10^{-j}\) with \(d_1,\ldots,d_m\) equal to \(s\), then \(x - a_s = \sum_{j > m} d_j 10^{-j} \in [0, 10^{-m}]\), so \(x \in I_s\). Conversely, if \(x \in I_s\), then \(10^{m}(x - a_s) \in [0,1]\), so we may write \(10^{m}(x-a_s) = \sum_{j=1}^{\infty} e_j 10^{-j}\) for some digits \(e_j\) (taking \(e_j = 9\) for all \(j\) in case \(10^m(x-a_s) = 1\)); then \(d_1,\ldots,d_m,e_1,e_2,\ldots\) is a decimal expansion of \(x\) beginning with \(s\).

(a). Let \(F\) denote the string of one hundred 4s. For each \(n \in \mathbb{Z}^{+}\) and each string \(p \in \{0,\ldots,9\}^{n-1}\), the set of \(x \in [0,1]\) having a decimal expansion whose digits in positions \(n, n+1, \ldots, n+99\) all equal 4 and whose first \(n-1\) digits are \(p\) is exactly \(I_{pF}\), where \(pF\) denotes the concatenation of \(p\) with \(F\). A number \(x\) lies in \(E\) exactly when some decimal expansion of \(x\) has one hundred consecutive 4s starting in some position \(n\), so

\begin{equation*} E = (0,1) \cap \bigcup_{n=1}^{\infty} \ \bigcup_{p \in \{0,\ldots,9\}^{n-1}} I_{pF}. \end{equation*}

For each \(n\) there are only \(10^{\,n-1}\) strings \(p\), so the double union is a countable union of closed intervals. Every closed interval is a Borel set and the collection of Borel sets is a \(\sigma\)-algebra, so \(E\) is the intersection of two Borel sets and hence is a Borel set.

(b). We claim \(|E| = 1\); equivalently \(|(0,1) \setminus E| = 0\).

Fix \(N \in \mathbb{Z}^{+}\). Call a string \(s \in \{0,\ldots,9\}^{100N}\) good if none of its \(N\) consecutive blocks of length 100 (positions \(100k+1,\ldots,100k+100\) for \(k = 0,\ldots,N-1\)) consists of one hundred 4s. The number of good strings is \((10^{100} - 1)^{N}\), since each block may be filled in \(10^{100}-1\) ways independently of the others. Let

\begin{equation*} U_N = \bigcup_{s \text{ good}} I_s . \end{equation*}

By subadditivity of outer measure (2.8) and \(|I_s| = 10^{-100N}\) (by 2.14),

\begin{equation*} |U_N| \le (10^{100}-1)^{N} \cdot 10^{-100N} = \Bigl(1 - 10^{-100}\Bigr)^{N}. \end{equation*}

Now suppose \(x \in (0,1) \setminus E\). Choose any decimal expansion \(d_1, d_2, \ldots\) of \(x\) and let \(s\) be the string of its first \(100N\) digits. If \(s\) were not good, then some block of one hundred consecutive positions of this expansion would consist entirely of 4s, so \(x\) would lie in \(E\); hence \(s\) is good. By the basic observation, \(x \in I_s \subseteq U_N\). Therefore \((0,1) \setminus E \subseteq U_N\), and so by 2.5

\begin{equation*} |(0,1) \setminus E| \le \bigl(1 - 10^{-100}\bigr)^{N} \end{equation*}

for every \(N \in \mathbb{Z}^{+}\). Letting \(N \to \infty\) gives \(|(0,1)\setminus E| = 0\).

Finally we compute \(|E|\) using only 2.4, 2.5, 2.8, and 2.14, so that no unproved value of \(|(0,1)|\) is needed. On one hand \(E \subseteq [0,1]\), so \(|E| \le |[0,1]| = 1\) by 2.5 and 2.14. On the other hand

\begin{equation*} [0,1] = E \cup \bigl( (0,1) \setminus E \bigr) \cup \{0,1\}, \end{equation*}

and \(|\{0,1\}| = 0\) by 2.4, so subadditivity (2.8) gives

\begin{equation*} 1 = |[0,1]| \le |E| + |(0,1)\setminus E| + |\{0,1\}| = |E| . \end{equation*}

Hence \(|E| = 1\): the set in (a) has Lebesgue measure \(1\).

Problem (2D.2)

Prove that there exists a bounded set \(A \subseteq \mathbb{R}\) such that \(|F| \le |A| - 1\) for every closed set \(F \subseteq A\).

Solution

We use the set \(V \subseteq [-1,1]\) constructed in the proof of 2.18: \(V\) contains exactly one element of each set \(\widetilde{a} = \{c \in [-1,1] : a - c \in \mathbb{Q}\}\), where \(a\) ranges over \([-1,1]\). Two properties established there will be used:

\begin{equation*} |V| > 0, \qquad \text{and} \qquad (r + V) \cap (r^{\prime} + V) = \varnothing \ \text{ for distinct } r, r^{\prime} \in \mathbb{Q}. \end{equation*}

(The disjointness was verified in the proof of 2.18 for the particular rationals used there, but the argument given uses only that \(r \ne r^{\prime}\) are rational: if \(t = r + v_1 = r^{\prime} + v_2\) with \(v_1, v_2 \in V\), then \(v_1 - v_2 = r^{\prime} - r \in \mathbb{Q}\), so \(v_1\) and \(v_2\) lie in the same set \(\widetilde{a}\), forcing \(v_1 = v_2\) and hence \(r = r^{\prime}\).)

Step 1: every closed subset of \(V\) has outer measure \(0\).

Suppose \(F \subseteq V\) is closed. Let \(r_1, r_2, \ldots\) be a sequence of distinct rational numbers in \([0,1]\). Because \(F \subseteq V\), the sets \(r_1 + F, r_2 + F, \ldots\) are disjoint. Each \(r_k + F\) is a closed subset of \(\mathbb{R}\) (a translate of a closed set is closed), hence a Borel set, and \(r_k + F \subseteq [-1, 2]\). Because outer measure is a measure on the Borel sets (2.68) and outer measure is translation invariant (2.7),

\begin{equation*} 3 = |[-1,2]| \ \ge\ \Bigl| \bigcup_{k=1}^{\infty} (r_k + F) \Bigr| = \sum_{k=1}^{\infty} |r_k + F| = \sum_{k=1}^{\infty} |F| , \end{equation*}

where the first inequality uses that outer measure preserves order (2.5) and 2.14. If \(|F| > 0\) the right side equals \(\infty\), a contradiction. Hence \(|F| = 0\).

Step 2: outer measure scales under dilation.

We claim \(|tB| = t|B|\) for every \(t > 0\) and every \(B \subseteq \mathbb{R}\), where \(tB = \{tb : b \in B\}\). If \(B \subseteq \bigcup_{k} (a_k, b_k)\), then \(tB \subseteq \bigcup_k (ta_k, tb_k)\) and \(\sum_k (tb_k - ta_k) = t \sum_k (b_k - a_k)\); taking the infimum over all such covers in the definition of outer measure (2.1) gives \(|tB| \le t|B|\). Applying this inequality with \(B\) replaced by \(tB\) and \(t\) replaced by \(1/t\) gives \(|B| \le \frac{1}{t}|tB|\), that is, \(t|B| \le |tB|\). Hence \(|tB| = t|B|\).

Step 3: the example.

Let \(c = |V|\), so \(0 < c < \infty\) (indeed \(c \le 2\) since \(V \subseteq [-1,1]\)), and set

\begin{equation*} A = \tfrac{1}{c} V = \Bigl\{ \tfrac{v}{c} : v \in V \Bigr\} . \end{equation*}

Then \(A \subseteq [-\frac{1}{c}, \frac{1}{c}]\), so \(A\) is bounded, and by Step 2,

\begin{equation*} |A| = \tfrac{1}{c}\,|V| = \tfrac{1}{c} \cdot c = 1 . \end{equation*}

Now suppose \(F \subseteq A\) is closed. Then \(cF\) is a closed subset of \(cA = V\), so \(|cF| = 0\) by Step 1, and therefore \(|F| = \frac{1}{c}|cF| = 0\) by Step 2. Thus

\begin{equation*} |F| = 0 = |A| - 1 \end{equation*}

for every closed set \(F \subseteq A\), which is the required inequality.

Problem (2D.3)

Prove that there exists a set \(A \subseteq \mathbb{R}\) such that \(|G \setminus A| = \infty\) for every open set \(G\) that contains \(A\).

Solution

We first record three general tools, then construct \(A\) from a Vitali set inside \((0,1)\).

Lemma 0 (outer measure of a bounded open interval). If \(a < b\) then \(|(a,b)| = b - a\).

Proof: By 2.5 and 2.14, \(|(a,b)| \le |[a,b]| = b-a\). For the reverse, let \(0 < \varepsilon < (b-a)/2\); then \([a+\varepsilon, b-\varepsilon] \subseteq (a,b)\), so by 2.5 and 2.14 again \(|(a,b)| \ge b - a - 2\varepsilon\). Letting \(\varepsilon \downarrow 0\) gives \(|(a,b)| \ge b-a\).

Lemma 1 (measurable hull). If \(S \subseteq \mathbb{R}\) and \(|S| < \infty\), then there is a Borel set \(B\) with \(S \subseteq B\) and \(|B| = |S|\).

Proof: For each \(k \in \mathbb{Z}^{+}\) the definition of outer measure (2.1) provides open intervals \(I_1, I_2, \ldots\) with \(S \subseteq \bigcup_j I_j\) and \(\sum_j \ell(I_j) < |S| + \frac{1}{k}\). Let \(U_k = \bigcup_j I_j\), an open set containing \(S\) with \(|U_k| \le \sum_j \ell(I_j) < |S| + \frac{1}{k}\). Put \(B = \bigcap_{k=1}^{\infty} U_k\). Then \(B\) is Borel, \(S \subseteq B\), and \(|S| \le |B| \le |U_k| < |S| + \frac{1}{k}\) for every \(k\), so \(|B| = |S|\).

Lemma 2 (splitting by a measurable set). If \(E \subseteq \mathbb{R}\) is Lebesgue measurable and \(S \subseteq \mathbb{R}\), then \(|S| = |S \cap E| + |S \setminus E|\).

Proof: The inequality \(\le\) is subadditivity (2.8). For \(\ge\) we may assume \(|S| < \infty\). Take \(B\) as in Lemma 1. Then \(S \cap E \subseteq B \cap E\) and \(S \setminus E \subseteq B \setminus E\), and \(B \cap E\), \(B \setminus E\) are disjoint Lebesgue measurable sets whose union is \(B\). Because outer measure is a measure on the Lebesgue measurable sets (2.72),

\begin{equation*} |S \cap E| + |S \setminus E| \le |B \cap E| + |B \setminus E| = |B| = |S| . \end{equation*}

Consequence. If \(E_1, \ldots, E_N\) are disjoint Lebesgue measurable sets and \(S \subseteq \mathbb{R}\), then \(|S| \ge \sum_{n=1}^{N} |S \cap E_n|\). This follows by induction on \(N\): Lemma 2 gives \(|S| = |S \cap E_1| + |S \setminus E_1|\), and \(S \setminus E_1\) contains \(S \cap E_n\) for \(n \ge 2\), so the induction hypothesis applies to \(S \setminus E_1\) and \(E_2, \ldots, E_N\).

A Vitali set in \((0,1)\). For \(a \in (0,1)\) let \(\widetilde{a} = \{c \in (0,1) : a - c \in \mathbb{Q}\}\). As in the proof of 2.18, distinct such sets are disjoint and they cover \((0,1)\); let \(V \subseteq (0,1)\) contain exactly one element of each distinct set \(\widetilde{a}\). Exactly as in 2.18, if \(r, r^{\prime}\) are distinct rational numbers then \((r+V) \cap (r^{\prime}+V) = \varnothing\).

If \(a \in (0,1)\) and \(v\) is the unique element of \(V \cap \widetilde{a}\), then \(a - v\) is a rational number in \((-1,1)\); hence, listing \(\mathbb{Q} \cap (-1,1)\) as \(q_1, q_2, \ldots\), we get \((0,1) \subseteq \bigcup_{k} (q_k + V)\). By Lemma 0, 2.5, 2.8, and 2.7,

\begin{equation*} 1 = |(0,1)| \le \sum_{k=1}^{\infty} |q_k + V| = \sum_{k=1}^{\infty} |V| , \end{equation*}

which forces \(|V| > 0\).

Claim 1: every closed set \(F \subseteq V\) has \(|F| = 0\). Let \(r_1, r_2, \ldots\) be distinct rationals in \([0,1]\). The sets \(r_k + F\) are disjoint closed (hence Borel) subsets of \((0,2)\), so by 2.68, 2.7, 2.5, and 2.14,

\begin{equation*} 2 \ge \Bigl| \bigcup_{k=1}^{\infty} (r_k + F) \Bigr| = \sum_{k=1}^{\infty} |r_k+F| = \sum_{k=1}^{\infty}|F| . \end{equation*}

Thus \(|F| = 0\).

Claim 2: \(V\) is not Lebesgue measurable. If it were, then since \(0 < |V| \le 1 < \infty\), condition (b) of 2.71 with \(\varepsilon = |V|/2\) would give a closed set \(F \subseteq V\) with \(|V \setminus F| < |V|/2\); then by 2.8 and Claim 1, \(|V| \le |F| + |V \setminus F| < 0 + |V|/2\), a contradiction.

Claim 3: \(c := |V| + |(0,1)\setminus V| - 1 > 0\). By subadditivity (2.8) and Lemma 0 we have \(|V| + |(0,1)\setminus V| \ge |(0,1)| = 1\), so \(c \ge 0\) (and \(c\) is a real number, since \(|V| \le 1\) and \(|(0,1)\setminus V| \le 1\) by 2.5 and Lemma 0). Suppose \(c = 0\), that is, \(|V| + |(0,1)\setminus V| = 1\). By Lemma 1 choose Borel sets \(B \supseteq V\) and \(C \supseteq (0,1)\setminus V\) with \(|B| = |V|\) and \(|C| = |(0,1)\setminus V|\); replacing \(B\) by \(B \cap (0,1)\) and \(C\) by \(C \cap (0,1)\) changes neither the inclusions nor the outer measures, so we may assume \(B, C \subseteq (0,1)\). Then \(B \cup C = (0,1)\). Since outer measure is a measure on the Borel sets (2.68),

\begin{equation*} |B| + |C| = |B| + |C \setminus B| + |C \cap B| = |B \cup C| + |B \cap C| = 1 + |B \cap C| , \end{equation*}

the last equality because \(B \cup C = (0,1)\) and \(|(0,1)| = 1\) by Lemma 0. (All the quantities here are at most \(1\), so the cancellations are legitimate.) As \(|B| + |C| = 1\), we get \(|B \cap C| = 0\). Now \((0,1) \setminus C \subseteq V\) because \((0,1)\setminus V \subseteq C\), hence

\begin{equation*} V = \bigl( (0,1) \setminus C \bigr) \cup (V \cap C), \end{equation*}

and \(V \cap C \subseteq B \cap C\) has outer measure \(0\) by 2.5, so \(V \cap C\) is Lebesgue measurable (in 2.70 take the Borel subset to be \(\varnothing\)). Since \((0,1)\setminus C\) is Borel and the Lebesgue measurable sets form a \(\sigma\)-algebra (2.72), \(V\) would be Lebesgue measurable, contradicting Claim 2. Hence \(c > 0\).

The construction. Let

\begin{equation*} A = \bigcup_{n=1}^{\infty} (n + V) . \end{equation*}

Since \(V \subseteq (0,1)\), we have \(n + V \subseteq (n, n+1)\), and therefore \(A \cap (n,n+1) = n+V\) for each \(n \in \mathbb{Z}^{+}\).

Let \(G\) be any open set with \(A \subseteq G\), and fix \(n \in \mathbb{Z}^{+}\). Put \(W = G \cap (n,n+1)\), an open set containing \(n+V\), and put \(E = (n,n+1) \setminus W\), a Borel set (hence Lebesgue measurable, by 2.70 with the Borel subset taken to be \(E\) itself) that is disjoint from \(n + V\) because \(n + V \subseteq W\).

First, applying Lemma 2 to \(S = (n+V) \cup E\) and the measurable set \(E\) gives

\begin{equation*} |E| + |n+V| = |S \cap E| + |S \setminus E| = |S| \le |(n,n+1)| = 1 , \end{equation*}

using \(S \cap E = E\) (as \(E \subseteq S\)), \(S \setminus E = n+V\) (as \(E\) and \(n+V\) are disjoint), \(S \subseteq (n,n+1)\), 2.5, and Lemma 0. Hence, by translation invariance (2.7), \(|E| \le 1 - |V|\).

Second, \((n,n+1) \setminus (n+V) \subseteq E \cup \bigl( W \setminus (n+V) \bigr)\), so by 2.8 and 2.7,

\begin{equation*} |(0,1) \setminus V| = \bigl| (n,n+1)\setminus (n+V) \bigr| \le |E| + \bigl| W \setminus (n+V) \bigr| . \end{equation*}

Combining the two displays,

\begin{equation*} \bigl| W \setminus (n+V) \bigr| \ge |(0,1)\setminus V| - (1 - |V|) = c . \end{equation*}

Because \(A \cap (n,n+1) = n+V\) and \(W = G \cap (n,n+1)\), we have \(W \setminus (n+V) = (G \setminus A) \cap (n,n+1)\). Hence

\begin{equation*} \bigl| (G\setminus A) \cap (n,n+1) \bigr| \ge c \qquad \text{for every } n \in \mathbb{Z}^{+} . \end{equation*}

Finally, the intervals \((1,2), (2,3), \ldots, (N, N+1)\) are disjoint Borel sets, so by the Consequence of Lemma 2,

\begin{equation*} |G \setminus A| \ge \sum_{n=1}^{N} \bigl| (G \setminus A) \cap (n,n+1) \bigr| \ge Nc \end{equation*}

for every \(N \in \mathbb{Z}^{+}\). Since \(c > 0\), letting \(N \to \infty\) gives \(|G \setminus A| = \infty\), as desired.

Problem (2D.4)

The phrase nontrivial interval is used to denote an interval of \(\mathbb{R}\) that contains more than one element. Recall that an interval might be open, closed, or neither.

(a) Prove that the union of each collection of nontrivial intervals of \(\mathbb{R}\) is the union of a countable subset of that collection.

(b) Prove that the union of each collection of nontrivial intervals of \(\mathbb{R}\) is a Borel set.

(c) Prove that there exists a collection of closed intervals of \(\mathbb{R}\) whose union is not a Borel set.

Solution

(a). Let \(\mathcal{C}\) be a collection of nontrivial intervals of \(\mathbb{R}\) and let \(U = \bigcup_{I \in \mathcal{C}} I\). For \(I \in \mathcal{C}\) let \(I^{\circ}\) denote the interior of \(I\); because \(I\) is an interval containing at least two points, \(I^{\circ}\) is a nonempty open interval and \(I \setminus I^{\circ}\) consists of at most the two endpoints of \(I\). Put

\begin{equation*} W = \bigcup_{I \in \mathcal{C}} I^{\circ} , \end{equation*}

an open subset of \(\mathbb{R}\) with \(W \subseteq U\).

Step 1: \(W\) is the union of countably many of the sets \(I^{\circ}\). Let \(D\) be the set of pairs \((p,q)\) of rational numbers with \(p < q\) such that \((p,q) \subseteq I^{\circ}\) for at least one \(I \in \mathcal{C}\); for each \((p,q) \in D\) choose one such interval and call it \(I_{p,q} \in \mathcal{C}\). The set \(D\) is countable, being a subset of \(\mathbb{Q} \times \mathbb{Q}\). If \(x \in W\), then \(x \in I^{\circ}\) for some \(I \in \mathcal{C}\), and since \(I^{\circ}\) is open there are rationals \(p < x < q\) with \((p,q) \subseteq I^{\circ}\); thus \((p,q) \in D\) and \(x \in (p,q) \subseteq I_{p,q}\). Hence

\begin{equation*} W \subseteq \bigcup_{(p,q) \in D} I_{p,q} \subseteq U . \end{equation*}

Step 2: \(U \setminus W\) is countable. Suppose \(x \in U \setminus W\). Then \(x \in I\) for some \(I \in \mathcal{C}\), and \(x \notin I^{\circ}\), so \(x\) is an endpoint of the nontrivial interval \(I\). Consequently \(I^{\circ}\) is a nonempty open interval having \(x\) as an endpoint, so there exists \(\delta > 0\) with \((x, x+\delta) \subseteq I^{\circ} \subseteq W\) or \((x-\delta, x) \subseteq I^{\circ} \subseteq W\). Let \(U_{+}\) be the set of \(x \in U \setminus W\) of the first kind and \(U_{-}\) the set of those of the second kind, so \(U \setminus W = U_{+} \cup U_{-}\).

For \(x \in U_{+}\) choose \(\delta_x > 0\) with \((x, x+\delta_x) \subseteq W\) and choose a rational \(r_x \in (x, x + \delta_x)\). The map \(x \mapsto r_x\) is injective on \(U_{+}\): if \(x, y \in U_{+}\) with \(x < y\) and \(r_x = r_y\), then \(x < y < r_y = r_x < x + \delta_x\), so \(y \in (x, x+\delta_x) \subseteq W\), contradicting \(y \notin W\). Hence \(U_{+}\) is countable, and symmetrically so is \(U_{-}\). Thus \(U \setminus W\) is countable.

Step 3. For each \(x \in U \setminus W\) choose \(J_x \in \mathcal{C}\) with \(x \in J_x\). Then

\begin{equation*} \mathcal{D} = \{ I_{p,q} : (p,q) \in D \} \cup \{ J_x : x \in U \setminus W \} \end{equation*}

is a countable subcollection of \(\mathcal{C}\), and

\begin{equation*} U = W \cup (U \setminus W) \subseteq \bigcup_{J \in \mathcal{D}} J \subseteq U . \end{equation*}

Therefore \(U = \bigcup_{J \in \mathcal{D}} J\), the union of a countable subcollection of \(\mathcal{C}\).

(b). Let \(\mathcal{C}\) and \(U\) be as above. By (a) there is a countable subcollection \(\mathcal{D} \subseteq \mathcal{C}\) with \(U = \bigcup_{J \in \mathcal{D}} J\). Every interval \(J\) of \(\mathbb{R}\) is a Borel set. Indeed, let \(a = \inf J \in [-\infty,\infty)\) and \(b = \sup J \in (-\infty,\infty]\) when \(J \ne \varnothing\). The set \(\{x \in \mathbb{R} : a < x < b\}\) is open, hence Borel, and \(J\) differs from it by at most the two points \(a\) and \(b\) (only those that are real and belong to \(J\)). Since every singleton \(\{a\} = \bigcap_{m=1}^{\infty} (a - \frac{1}{m}, a + \frac{1}{m})\) is a countable intersection of open sets and hence Borel, and since the Borel sets form a \(\sigma\)-algebra, \(J\) is Borel; and \(\varnothing\) is Borel. Therefore the countable union \(U = \bigcup_{J \in \mathcal{D}} J\) is a Borel set.

(c). Note first that part (c) drops the word nontrivial: for \(x \in \mathbb{R}\) the singleton \(\{x\} = [x,x]\) is a closed interval, a trivial one. By 2.67 there exists a set \(S \subseteq \mathbb{R}\) with \(|S| < \infty\) that is not a Borel set.

Now let

\begin{equation*} \mathcal{C} = \bigl\{ [x,x] : x \in S \bigr\} , \end{equation*}

a collection of closed intervals of \(\mathbb{R}\). Its union is

\begin{equation*} \bigcup_{x \in S} \{x\} = S , \end{equation*}

which is not a Borel set. This shows that the hypothesis of nontriviality cannot be dropped in (a) or (b): the singletons \(\{x\}\) have empty interior, which is exactly what makes the counting argument of (a) break down.

Problem (2D.5)

Prove that if \(A \subseteq \mathbb{R}\) is Lebesgue measurable, then there exists an increasing sequence \(F_1 \subseteq F_2 \subseteq \cdots\) of closed sets contained in \(A\) such that

\begin{equation*} \Bigl| A \setminus \bigcup_{k=1}^{\infty} F_k \Bigr| = 0 . \end{equation*}

Solution

Suppose \(A \subseteq \mathbb{R}\) is Lebesgue measurable. Because (a) implies (b) in 2.71, for each \(k \in \mathbb{Z}^{+}\) there exists a closed set \(E_k \subseteq A\) with

\begin{equation*} |A \setminus E_k| < \tfrac{1}{k} . \end{equation*}

The sets \(E_k\) produced this way need not be increasing, so define

\begin{equation*} F_k = E_1 \cup E_2 \cup \cdots \cup E_k \qquad \text{for } k \in \mathbb{Z}^{+} . \end{equation*}

Each \(F_k\) is a union of finitely many closed sets and hence is closed. Each \(F_k\) is contained in \(A\) because each \(E_j \subseteq A\). Clearly \(F_1 \subseteq F_2 \subseteq \cdots\), and

\begin{equation*} \bigcup_{k=1}^{\infty} F_k = \bigcup_{k=1}^{\infty} E_k . \end{equation*}

Now fix \(n \in \mathbb{Z}^{+}\). Because \(E_n \subseteq \bigcup_{k=1}^{\infty} F_k\), we have

\begin{equation*} A \setminus \bigcup_{k=1}^{\infty} F_k \subseteq A \setminus E_n , \end{equation*}

and therefore, because outer measure preserves order (2.5),

\begin{equation*} \Bigl| A \setminus \bigcup_{k=1}^{\infty} F_k \Bigr| \le |A \setminus E_n| < \tfrac{1}{n} . \end{equation*}

This holds for every \(n \in \mathbb{Z}^{+}\), so

\begin{equation*} \Bigl| A \setminus \bigcup_{k=1}^{\infty} F_k \Bigr| = 0 , \end{equation*}

which is what we wanted. (This is exactly the argument used to prove that (b) implies (c) in 2.71, with the additional observation that replacing the closed sets by their partial unions makes the sequence increasing without changing the union.)

Problem (2D.6)

Suppose \(A \subseteq \mathbb{R}\) and \(|A| < \infty\). Prove that \(A\) is Lebesgue measurable if and only if for every \(\varepsilon > 0\) there exists a set \(G\) that is the union of finitely many disjoint bounded open intervals such that \(|A \setminus G| + |G \setminus A| < \varepsilon\).

Solution

Suppose \(A\) is Lebesgue measurable, and let \(\varepsilon > 0\).

Because (a) implies (e) in 2.71, there is an open set \(U \supseteq A\) with

\begin{equation*} |U \setminus A| < \tfrac{\varepsilon}{2} . \end{equation*}

By subadditivity of outer measure (2.8), \(|U| \le |A| + |U \setminus A| < \infty\).

Every nonempty open subset of \(\mathbb{R}\) is the union of a countable disjoint collection of nonempty open intervals, namely its connected components; this is the fact invoked in the proof of 2.63, where an arbitrary open subset \(W\) of \(\mathbb{R}\) is written as \(W = \bigcup_{n=1}^{\infty} I_n\) for a sequence of disjoint open intervals. Dispose first of the degenerate case: if \(U = \varnothing\) then \(A \subseteq U\) forces \(A = \varnothing\), and \(G = (0, \varepsilon/3)\) is a union of one bounded open interval with \(|A \setminus G| + |G \setminus A| = |G| \le |[0,\varepsilon/3]| = \varepsilon/3 < \varepsilon\) by 2.5 and 2.14. So assume \(U \ne \varnothing\) and write \(U = \bigcup_{n} I_n\) with \(I_1, I_2, \ldots\) a finite or infinite sequence of disjoint nonempty open intervals; pad a finite list with copies of \(\varnothing\) so that the index set is \(\mathbb{Z}^{+}\).

Each \(I_n\) is bounded: an unbounded nonempty open interval contains closed intervals of arbitrarily large length, so 2.5 and 2.14 would give \(|U| \ge |I_n| = \infty\), contradicting \(|U| < \infty\). Because the \(I_n\) are disjoint Borel sets and outer measure is a measure on the Borel sets (2.68),

\begin{equation*} \sum_{n=1}^{\infty} |I_n| = |U| < \infty . \end{equation*}

Since this series converges, choose \(N \in \mathbb{Z}^{+}\) with

\begin{equation*} \sum_{n > N} |I_n| < \tfrac{\varepsilon}{2}, \end{equation*}

and set \(G = I_1 \cup \cdots \cup I_N\), a union of finitely many disjoint bounded open intervals. (Discarding any empty ones changes nothing.) Then:

  • \(G \setminus A \subseteq U \setminus A\), so \(|G \setminus A| \le |U \setminus A| < \frac{\varepsilon}{2}\) by 2.5.

  • Since \(A \subseteq U\), we have \(A \setminus G \subseteq U \setminus G = \bigcup_{n > N} I_n\), so by 2.5 and 2.8, \(|A \setminus G| \le \sum_{n>N} |I_n| < \frac{\varepsilon}{2}\).

Adding gives \(|A \setminus G| + |G \setminus A| < \varepsilon\), as required. (Note that this direction did not use \(|A| < \infty\) beyond guaranteeing \(|U| < \infty\), which is where the hypothesis is needed: without it the tail of \(\sum |I_n|\) need not be small.)

Conversely, suppose that for every \(\varepsilon > 0\) such a set \(G\) exists.

Let \(\varepsilon > 0\), and choose \(G\), a union of finitely many disjoint bounded open intervals, with

\begin{equation*} |A \setminus G| + |G \setminus A| < \tfrac{\varepsilon}{2} . \end{equation*}

In particular \(|A \setminus G| < \frac{\varepsilon}{2}\). By the definition of outer measure (2.1) there exist open intervals \(J_1, J_2, \ldots\) with

\begin{equation*} A \setminus G \subseteq \bigcup_{k} J_k \quad \text{and} \quad \sum_{k} \ell(J_k) < \tfrac{\varepsilon}{2} . \end{equation*}

Let \(H = \bigcup_k J_k\). Then \(H\) is open, \(A \setminus G \subseteq H\), and \(|H| \le \sum_k \ell(J_k) < \frac{\varepsilon}{2}\) by 2.8 and 2.14.

Set \(V = G \cup H\). Then \(V\) is open, and

\begin{equation*} A = (A \cap G) \cup (A \setminus G) \subseteq G \cup H = V , \end{equation*}

so \(V \supseteq A\). Moreover

\begin{equation*} V \setminus A = (G \cup H) \setminus A \subseteq (G \setminus A) \cup H , \end{equation*}

so by 2.5 and 2.8,

\begin{equation*} |V \setminus A| \le |G \setminus A| + |H| < \tfrac{\varepsilon}{2} + \tfrac{\varepsilon}{2} = \varepsilon . \end{equation*}

Thus for every \(\varepsilon > 0\) there is an open set \(V \supseteq A\) with \(|V \setminus A| < \varepsilon\); that is, \(A\) satisfies condition (e) of 2.71. Because (e) implies (a) in 2.71, the set \(A\) is Lebesgue measurable.

Problem (2D.7)

Prove that if \(A \subseteq \mathbb{R}\) is Lebesgue measurable, then there exists a decreasing sequence \(G_1 \supseteq G_2 \supseteq \cdots\) of open sets containing \(A\) such that

\begin{equation*} \Bigl| \Bigl( \bigcap_{k=1}^{\infty} G_k \Bigr) \setminus A \Bigr| = 0 . \end{equation*}

Solution

Suppose \(A \subseteq \mathbb{R}\) is Lebesgue measurable. Because (a) implies (e) in 2.71, for each \(k \in \mathbb{Z}^{+}\) there exists an open set \(U_k \supseteq A\) with

\begin{equation*} |U_k \setminus A| < \tfrac{1}{k} . \end{equation*}

These sets need not be decreasing, so define

\begin{equation*} G_k = U_1 \cap U_2 \cap \cdots \cap U_k \qquad \text{for } k \in \mathbb{Z}^{+} . \end{equation*}

Each \(G_k\) is an intersection of finitely many open sets and hence is open. Each \(G_k\) contains \(A\) because every \(U_j\) contains \(A\). Clearly \(G_1 \supseteq G_2 \supseteq \cdots\), and

\begin{equation*} \bigcap_{k=1}^{\infty} G_k = \bigcap_{k=1}^{\infty} U_k . \end{equation*}

Now fix \(n \in \mathbb{Z}^{+}\). Because \(\bigcap_{k=1}^{\infty} G_k \subseteq U_n\), we have

\begin{equation*} \Bigl( \bigcap_{k=1}^{\infty} G_k \Bigr) \setminus A \subseteq U_n \setminus A , \end{equation*}

and therefore, because outer measure preserves order (2.5),

\begin{equation*} \Bigl| \Bigl( \bigcap_{k=1}^{\infty} G_k \Bigr) \setminus A \Bigr| \le |U_n \setminus A| < \tfrac{1}{n} . \end{equation*}

Since this holds for every \(n \in \mathbb{Z}^{+}\),

\begin{equation*} \Bigl| \Bigl( \bigcap_{k=1}^{\infty} G_k \Bigr) \setminus A \Bigr| = 0 , \end{equation*}

as desired. (This is the argument proving that (e) implies (f) in 2.71, improved by replacing the open sets with their partial intersections, which makes the sequence decreasing without changing the intersection.)

Problem (2D.8)

Prove that the collection of Lebesgue measurable subsets of \(\mathbb{R}\) is translation invariant. More precisely, prove that if \(A \subseteq \mathbb{R}\) is Lebesgue measurable and \(t \in \mathbb{R}\), then \(t + A\) is Lebesgue measurable.

Solution

Recall the definition 2.70: a set \(A \subseteq \mathbb{R}\) is Lebesgue measurable if there exists a Borel set \(B \subseteq A\) such that \(|A \setminus B| = 0\). So we need two ingredients: that the Borel sets are translation invariant, and that outer measure is translation invariant (the latter is 2.7).

Step 1: the collection of Borel sets is translation invariant.

Fix \(t \in \mathbb{R}\) and let

\begin{equation*} \mathcal{S} = \{E \subseteq \mathbb{R} : t + E \text{ is a Borel set}\}. \end{equation*}

We claim \(\mathcal{S}\) is a \(\sigma\)-algebra on \(\mathbb{R}\).

  • \(\emptyset \in \mathcal{S}\), because \(t + \emptyset = \emptyset\), which is a Borel set.
  • Suppose \(E \in \mathcal{S}\). Because \(x \mapsto t + x\) is a bijection of \(\mathbb{R}\) onto \(\mathbb{R}\), we have \(t + (\mathbb{R} \setminus E) = \mathbb{R} \setminus (t + E)\). The right-hand side is the complement of a Borel set and hence is a Borel set. Thus \(\mathbb{R} \setminus E \in \mathcal{S}\).
  • Suppose \(E_1, E_2, \ldots \in \mathcal{S}\). Then

\begin{equation*} t + \bigcup_{k=1}^{\infty} E_k = \bigcup_{k=1}^{\infty} (t + E_k), \end{equation*}

which is a countable union of Borel sets and hence is a Borel set. Thus \(\bigcup_{k=1}^{\infty} E_k \in \mathcal{S}\).

Hence \(\mathcal{S}\) is a \(\sigma\)-algebra on \(\mathbb{R}\). Furthermore, if \(G \subseteq \mathbb{R}\) is open then \(t + G\) is open (translation is a homeomorphism of \(\mathbb{R}\); concretely, if \(x \in t + G\) then \(x - t \in G\), so \((x - t - \delta, x - t + \delta) \subseteq G\) for some \(\delta > 0\), and then \((x - \delta, x + \delta) \subseteq t + G\)), and every open set is a Borel set. Thus \(\mathcal{S}\) contains all open subsets of \(\mathbb{R}\).

Because the collection of Borel sets is by definition (2.29) the smallest \(\sigma\)-algebra on \(\mathbb{R}\) containing all the open sets, and \(\mathcal{S}\) is a \(\sigma\)-algebra containing all the open sets, we conclude that every Borel set belongs to \(\mathcal{S}\). In other words, if \(B\) is a Borel set and \(t \in \mathbb{R}\), then \(t + B\) is a Borel set.

Step 2: conclusion.

Suppose \(A \subseteq \mathbb{R}\) is Lebesgue measurable and \(t \in \mathbb{R}\). By 2.70 there exists a Borel set \(B \subseteq A\) with \(|A \setminus B| = 0\).

By Step 1, \(t + B\) is a Borel set. Because \(B \subseteq A\) we have \(t + B \subseteq t + A\).

Because \(x \mapsto t + x\) is injective,

\begin{equation*} (t + A) \setminus (t + B) = t + (A \setminus B). \end{equation*}

Indeed, \(x \in (t+A) \setminus (t+B)\) if and only if \(x - t \in A\) and \(x - t \notin B\), i.e. if and only if \(x - t \in A \setminus B\).

Hence, by the translation invariance of outer measure (2.7),

\begin{equation*} \bigl| (t + A) \setminus (t + B) \bigr| = \bigl| t + (A \setminus B) \bigr| = |A \setminus B| = 0 . \end{equation*}

Thus \(t + B\) is a Borel set contained in \(t + A\) whose complement in \(t + A\) has outer measure \(0\). By 2.70, \(t + A\) is Lebesgue measurable.

Problem (2D.9)

Prove that the collection of Lebesgue measurable subsets of \(\mathbb{R}\) is dilation invariant. More precisely, prove that if \(A \subseteq \mathbb{R}\) is Lebesgue measurable and \(t \in \mathbb{R}\), then \(tA\) (which is defined to be \(\{ta : a \in A\}\)) is Lebesgue measurable.

Solution

The proof parallels the one for translation (Exercise 8 in this section), using in place of 2.7 the dilation formula \(|tA| = |t| \, |A|\) from Exercise 2 in Section 2A (with the convention \(0 \cdot \infty = 0\)).

The degenerate case \(t = 0\). If \(A = \emptyset\) then \(0A = \emptyset\); otherwise \(0A = \{0\}\). In either case \(0A\) is a closed set, hence a Borel set, hence Lebesgue measurable. (Alternatively: \(|0A| = 0\) by 2.3, and every set of outer measure \(0\) is Lebesgue measurable, taking \(B = \emptyset\) in 2.70.) So assume from now on that \(t \neq 0\).

Step 1: the collection of Borel sets is dilation invariant.

Fix \(t \in \mathbb{R}\) with \(t \neq 0\) and let

\begin{equation*} \mathcal{S} = \{E \subseteq \mathbb{R} : tE \text{ is a Borel set}\}. \end{equation*}

We claim \(\mathcal{S}\) is a \(\sigma\)-algebra on \(\mathbb{R}\). Note first that because \(t \neq 0\), the map \(x \mapsto tx\) is a bijection of \(\mathbb{R}\) onto \(\mathbb{R}\), with inverse \(x \mapsto \frac{1}{t}x\).

  • \(\emptyset \in \mathcal{S}\), because \(t\emptyset = \emptyset\) is a Borel set.
  • Suppose \(E \in \mathcal{S}\). By the bijectivity just noted, \(t(\mathbb{R} \setminus E) = \mathbb{R} \setminus tE\), which is a Borel set. Thus \(\mathbb{R} \setminus E \in \mathcal{S}\).
  • Suppose \(E_1, E_2, \ldots \in \mathcal{S}\). Then

\begin{equation*} t \bigcup_{k=1}^{\infty} E_k = \bigcup_{k=1}^{\infty} tE_k , \end{equation*}

a countable union of Borel sets, hence a Borel set. Thus \(\bigcup_{k=1}^{\infty} E_k \in \mathcal{S}\).

Hence \(\mathcal{S}\) is a \(\sigma\)-algebra on \(\mathbb{R}\). If \(G \subseteq \mathbb{R}\) is open, then \(tG\) is open: if \(x \in tG\) then \(\frac{x}{t} \in G\), so \(\bigl(\frac{x}{t} - \delta, \frac{x}{t} + \delta\bigr) \subseteq G\) for some \(\delta > 0\), and then the interval of radius \(|t|\delta\) centered at \(x\) is contained in \(tG\). Thus \(\mathcal{S}\) contains every open subset of \(\mathbb{R}\).

Because the collection of Borel sets is the smallest \(\sigma\)-algebra on \(\mathbb{R}\) containing the open sets (2.29), we conclude that every Borel set is in \(\mathcal{S}\). In other words, if \(B\) is a Borel set and \(t \neq 0\), then \(tB\) is a Borel set.

Step 2: conclusion.

Suppose \(A \subseteq \mathbb{R}\) is Lebesgue measurable and \(t \neq 0\). By 2.70 there is a Borel set \(B \subseteq A\) with \(|A \setminus B| = 0\).

By Step 1, \(tB\) is a Borel set, and clearly \(tB \subseteq tA\). Because \(x \mapsto tx\) is injective,

\begin{equation*} tA \setminus tB = t(A \setminus B); \end{equation*}

indeed \(x \in tA \setminus tB\) if and only if \(\frac{x}{t} \in A\) and \(\frac{x}{t} \notin B\), i.e. if and only if \(\frac{x}{t} \in A \setminus B\).

By Exercise 2 in Section 2A,

\begin{equation*} |tA \setminus tB| = |t(A \setminus B)| = |t| \, |A \setminus B| = |t| \cdot 0 = 0 . \end{equation*}

Thus \(tB\) is a Borel set contained in \(tA\) with \(|tA \setminus tB| = 0\), so \(tA\) is Lebesgue measurable by 2.70.

Remark on the cited dilation formula. For completeness, here is the short proof of \(|tA| = |t| \, |A|\) for \(t \neq 0\). If \(I_1, I_2, \ldots\) are open intervals whose union contains \(A\), then \(tI_1, tI_2, \ldots\) are open intervals whose union contains \(tA\), and \(\ell(tI_k) = |t| \, \ell(I_k)\). Hence \(|tA| \leq \sum_{k=1}^{\infty} |t| \, \ell(I_k) = |t| \sum_{k=1}^{\infty} \ell(I_k)\); taking the infimum over all such covers gives \(|tA| \leq |t| \, |A|\). Applying this inequality to \(\frac{1}{t}\) and the set \(tA\) gives \(|A| = \bigl|\frac{1}{t}(tA)\bigr| \leq \frac{1}{|t|} |tA|\), i.e. \(|t| \, |A| \leq |tA|\). Combining the two inequalities gives equality.

Problem (2D.10)

Prove that if \(A\) and \(B\) are disjoint subsets of \(\mathbb{R}\) and \(B\) is Lebesgue measurable, then \(|A \cup B| = |A| + |B|\).

Solution

Note that no measurability is assumed of \(A\); the point of the exercise is that a single measurable set already splits outer measure additively off any disjoint companion.

The inequality \(|A \cup B| \leq |A| + |B|\) is immediate from countable subadditivity of outer measure (2.8), applied to the sequence \(A, B, \emptyset, \emptyset, \ldots\). So the whole content is the reverse inequality.

A lemma: every set has a Borel hull.

We first record the following, which we use again in Exercise 12 of this section.

Lemma. If \(S \subseteq \mathbb{R}\), then there exists a Borel set \(E\) with \(S \subseteq E\) and \(|E| = |S|\).

Proof of Lemma. If \(|S| = \infty\), take \(E = \mathbb{R}\); then \(S \subseteq E\) and \(|E| = \infty = |S|\), since \(|S| \leq |E|\) by 2.5. So suppose \(|S| < \infty\) and let \(n \in \mathbb{Z}^+\). Because \(|S|\) is defined in 2.2 as an infimum, there exist open intervals \(I_1, I_2, \ldots\) with \(S \subseteq \bigcup_{k=1}^{\infty} I_k\) and

\begin{equation*} \sum_{k=1}^{\infty} \ell(I_k) < |S| + \tfrac{1}{n}. \end{equation*}

Let \(G_n = \bigcup_{k=1}^{\infty} I_k\), an open set containing \(S\). Because \(|I| \leq \ell(I)\) for every open interval \(I\) (take the cover \(I, \emptyset, \emptyset, \ldots\) in 2.2), countable subadditivity (2.8) gives

\begin{equation*} |G_n| \leq \sum_{k=1}^{\infty} |I_k| \leq \sum_{k=1}^{\infty} \ell(I_k) < |S| + \tfrac{1}{n}. \end{equation*}

Now let \(E = \bigcap_{n=1}^{\infty} G_n\). Each \(G_n\) is open, hence Borel, so \(E\) is Borel; and \(S \subseteq E\) because \(S \subseteq G_n\) for every \(n\). By 2.5 we have \(|E| \leq |G_n| < |S| + \frac{1}{n}\) for every \(n \in \mathbb{Z}^+\), hence \(|E| \leq |S|\). The reverse inequality \(|S| \leq |E|\) holds by 2.5. Thus \(|E| = |S|\), proving the Lemma.

Proof of the exercise.

Suppose \(A, B \subseteq \mathbb{R}\) are disjoint and \(B\) is Lebesgue measurable.

Because \(B\) is Lebesgue measurable, 2.70 provides a Borel set \(D \subseteq B\) with \(|B \setminus D| = 0\). Then

\begin{equation*} |B| = |D \cup (B \setminus D)| \leq |D| + |B \setminus D| = |D| \leq |B|, \end{equation*}

using 2.8 for the first inequality and 2.5 for the last. Hence \(|D| = |B|\).

By the Lemma, choose a Borel set \(E\) with \(A \cup B \subseteq E\) and \(|E| = |A \cup B|\).

Since \(D \subseteq B \subseteq E\), the sets \(D\) and \(E \setminus D\) are disjoint Borel sets whose union is \(E\). Because outer measure is a measure on the Borel sets (2.68), it is in particular finitely additive there, so

\begin{equation*} |E| = |D| + |E \setminus D|. \end{equation*}

Now observe that \(A \subseteq E \setminus D\): indeed \(A \subseteq A \cup B \subseteq E\), and \(A \cap D \subseteq A \cap B = \emptyset\) because \(D \subseteq B\) and \(A, B\) are disjoint. Hence \(|A| \leq |E \setminus D|\) by 2.5.

Putting these together,

\begin{equation*} |A \cup B| = |E| = |D| + |E \setminus D| \geq |B| + |A| . \end{equation*}

Combined with the subadditivity inequality \(|A \cup B| \leq |A| + |B|\) noted at the start, this gives

\begin{equation*} |A \cup B| = |A| + |B| , \end{equation*}

as desired. (All the arithmetic above takes place in \([0, \infty]\), so no finiteness assumptions are needed.)

Problem (2D.11)

Prove that if \(A \subseteq \mathbb{R}\) and \(|A| > 0\), then there exists a subset of \(A\) that is not Lebesgue measurable.

Solution

This strengthens the existence of a nonmeasurable set (which follows from 2.18 and 2.72): every set of positive outer measure, no matter how small or how strangely placed, contains a nonmeasurable subset. As in 2.18, the Axiom of Choice is used.

A Vitali set with representatives in \([0,1]\).

Define a relation on \(\mathbb{R}\) by declaring \(x\) and \(y\) equivalent when \(x - y \in \mathbb{Q}\). This is an equivalence relation (reflexivity, symmetry, and transitivity all follow from the fact that \(\mathbb{Q}\) is closed under negation and addition). Write \(\tilde{x} = x + \mathbb{Q}\) for the equivalence class of \(x\); distinct classes are disjoint.

Each class \(\tilde{x}\) meets \([0,1]\): choose a rational \(q\) with \(-x \leq q \leq -x + 1\), which exists because \(\mathbb{Q}\) is dense in \(\mathbb{R}\); then \(x + q \in \tilde{x} \cap [0,1]\).

The sets \(\tilde{x} \cap [0,1]\), as \(\tilde{x}\) ranges over the distinct equivalence classes, are disjoint and (by the previous paragraph) nonempty. So the Axiom of Choice provides a set \(V\) containing exactly one element of \(\tilde{x} \cap [0,1]\) for each class \(\tilde{x}\). Then \(V \subseteq [0,1]\), and because \(V \cap \tilde{x} = V \cap (\tilde{x} \cap [0,1])\), the set \(V\) contains exactly one element of each equivalence class.

Two properties of \(V\):

  • (i) \(\mathbb{R} = \bigcup_{r \in \mathbb{Q}} (r + V)\). Indeed, given \(x \in \mathbb{R}\), let \(v\) be the unique element of \(V \cap \tilde{x}\). Then \(x - v \in \mathbb{Q}\) and \(x = (x - v) + v \in (x - v) + V\).
  • (ii) The sets \(r + V\), for \(r \in \mathbb{Q}\), are pairwise disjoint. Indeed, if \(t \in (r + V) \cap (s + V)\) with \(r, s \in \mathbb{Q}\), write \(t = r + v_1 = s + v_2\) with \(v_1, v_2 \in V\). Then \(v_1 - v_2 = s - r \in \mathbb{Q}\), so \(v_1\) and \(v_2\) lie in the same equivalence class; because \(V\) contains exactly one element of each class, \(v_1 = v_2\), and therefore \(r = s\).

Proof of the exercise.

Suppose \(A \subseteq \mathbb{R}\) with \(|A| > 0\), and suppose toward a contradiction that every subset of \(A\) is Lebesgue measurable.

For \(r \in \mathbb{Q}\) let

\begin{equation*} A_r = A \cap (r + V). \end{equation*}

Each \(A_r\) is a subset of \(A\), hence is Lebesgue measurable by our assumption. By (i) we have \(A = \bigcup_{r \in \mathbb{Q}} A_r\), and by (ii) this union is disjoint. The set \(\mathbb{Q}\) is countable, so this is a countable disjoint union of Lebesgue measurable sets. Because outer measure is a measure on the \(\sigma\)-algebra of Lebesgue measurable sets (2.72), it is countably additive there, so

\begin{equation*} 0 < |A| = \sum_{r \in \mathbb{Q}} |A_r| . \end{equation*}

Hence there exists \(r \in \mathbb{Q}\) with \(|A_r| > 0\). Fix such an \(r\) and set

\begin{equation*} W = (-r) + A_r . \end{equation*}

By Exercise 8 in this section, \(W\) is Lebesgue measurable, and by the translation invariance of outer measure (2.7), \(|W| = |A_r| > 0\). Moreover \(A_r \subseteq r + V\) gives \(W \subseteq V \subseteq [0,1]\).

Now let \(q_1, q_2, \ldots\) be a sequence of distinct rational numbers in \([0,1]\) (for instance an enumeration of \(\mathbb{Q} \cap [0,1]\)). For each \(k \in \mathbb{Z}^+\):

  • \(q_k + W\) is Lebesgue measurable (Exercise 8) and \(|q_k + W| = |W|\) (2.7);
  • \(q_k + W \subseteq q_k + V\), so by (ii) the sets \(q_1 + W, q_2 + W, \ldots\) are pairwise disjoint;
  • \(q_k + W \subseteq [0,1] + [0,1] \subseteq [0,2]\), since \(W \subseteq [0,1]\) and \(q_k \in [0,1]\).

Therefore \(\bigcup_{k=1}^{\infty} (q_k + W) \subseteq [0,2]\), and by countable additivity of outer measure on the Lebesgue measurable sets (2.72), together with 2.5 and 2.14,

\begin{equation*} \sum_{k=1}^{\infty} |W| \;=\; \sum_{k=1}^{\infty} |q_k + W| \;=\; \Bigl| \bigcup_{k=1}^{\infty} (q_k + W) \Bigr| \;\leq\; |[0,2]| \;=\; 2 . \end{equation*}

But \(|W| > 0\), so the left side of the display above equals \(\infty\), a contradiction.

Hence our assumption was false: some subset of \(A\) is not Lebesgue measurable.

Problem (2D.12)

Suppose \(b < c\) and \(A \subseteq (b, c)\). Prove that \(A\) is Lebesgue measurable if and only if \(|A| + |(b,c) \setminus A| = c - b\).

Solution

Throughout we use that \(|(b,c)| = c - b\) (Exercise 6 in Section 2A). For completeness: \((b,c) \subseteq [b,c]\) gives \(|(b,c)| \leq |[b,c]| = c - b\) by 2.5 and 2.14, while \([b,c] = (b,c) \cup \{b,c\}\) gives \(c - b = |[b,c]| \leq |(b,c)| + |\{b,c\}| = |(b,c)|\) by 2.8 and 2.3.

Suppose \(A\) is Lebesgue measurable.

The sets \((b,c) \setminus A\) and \(A\) are disjoint, and \(A\) is Lebesgue measurable. By Exercise 10 in this section (applied with the measurable set in the role of \(B\)),

\begin{equation*} \bigl| \bigl( (b,c) \setminus A \bigr) \cup A \bigr| = |(b,c) \setminus A| + |A| . \end{equation*}

Because \(A \subseteq (b,c)\), the union on the left equals \((b,c)\). Hence

\begin{equation*} |A| + |(b,c) \setminus A| = |(b,c)| = c - b . \end{equation*}

Suppose \(|A| + |(b,c) \setminus A| = c - b\).

Both \(A\) and \((b,c) \setminus A\) are subsets of \((b,c)\), so by 2.5 each of \(|A|\) and \(|(b,c) \setminus A|\) is at most \(c - b\); in particular both are finite, so the arithmetic below is legitimate.

By the Lemma proved in Exercise 10 of this section, there is a Borel set \(E_0\) with \((b,c) \setminus A \subseteq E_0\) and \(|E_0| = |(b,c) \setminus A|\). Replace \(E_0\) by

\begin{equation*} E = E_0 \cap (b,c), \end{equation*}

which is again Borel and still contains \((b,c) \setminus A\); and since

\begin{equation*} |(b,c) \setminus A| \leq |E| \leq |E_0| = |(b,c) \setminus A| \end{equation*}

by 2.5, we still have \(|E| = |(b,c) \setminus A|\). Thus \(E\) is a Borel subset of \((b,c)\) containing \((b,c) \setminus A\) with \(|E| = |(b,c) \setminus A|\).

Now let

\begin{equation*} B = (b,c) \setminus E . \end{equation*}

Then \(B\) is a Borel set. Moreover \(B \subseteq A\): if \(x \in B\) then \(x \in (b,c)\) and \(x \notin E\); since \((b,c) \setminus A \subseteq E\), we get \(x \notin (b,c) \setminus A\), and hence \(x \in A\).

The sets \(B\) and \(E\) are disjoint Borel sets whose union is \((b,c)\). Because outer measure is a measure on the Borel sets (2.68), it is finitely additive there, so

\begin{equation*} |B| + |E| = |(b,c)| = c - b . \end{equation*}

Since \(|E| = |(b,c) \setminus A|\) is finite, we may subtract and use the hypothesis:

\begin{equation*} |B| = c - b - |(b,c) \setminus A| = |A| . \end{equation*}

Finally, \(A \setminus B\) and \(B\) are disjoint and \(B\) is Lebesgue measurable (it is a Borel set). By Exercise 10 in this section,

\begin{equation*} \bigl| (A \setminus B) \cup B \bigr| = |A \setminus B| + |B| . \end{equation*}

Because \(B \subseteq A\), the union on the left equals \(A\). Hence

\begin{equation*} |A| = |A \setminus B| + |B| = |A \setminus B| + |A| . \end{equation*}

Since \(|A| \leq c - b < \infty\), we may cancel \(|A|\) to conclude \(|A \setminus B| = 0\).

Thus \(B\) is a Borel set with \(B \subseteq A\) and \(|A \setminus B| = 0\), which by definition 2.70 says exactly that \(A\) is Lebesgue measurable.

Problem (2D.13)

Suppose \(A \subseteq \mathbb{R}\). Prove that \(A\) is Lebesgue measurable if and only if

\begin{equation*} |(-n,n) \cap A| + |(-n,n) \setminus A| = 2n \end{equation*}

for every \(n \in \mathbb{Z}^+\).

Solution

For each \(n \in \mathbb{Z}^+\) set

\begin{equation*} A_n = (-n,n) \cap A . \end{equation*}

Then \(A_n \subseteq (-n,n)\), and

\begin{equation*} (-n,n) \setminus A_n = (-n,n) \setminus A , \end{equation*}

because removing \(A\) from \((-n,n)\) is the same as removing the part of \(A\) that lies inside \((-n,n)\). So the displayed condition in the exercise says precisely

\begin{equation*} |A_n| + |(-n,n) \setminus A_n| = 2n = n - (-n) \end{equation*}

for every \(n \in \mathbb{Z}^+\). Exercise 12 in this section (with \(b = -n\) and \(c = n\)) says that this condition holds if and only if \(A_n\) is Lebesgue measurable. So the exercise reduces to showing:

\begin{equation*} A \text{ is Lebesgue measurable} \iff A_n \text{ is Lebesgue measurable for every } n \in \mathbb{Z}^+ . \end{equation*}

Suppose \(A\) is Lebesgue measurable.

The collection \(\mathcal{L}\) of Lebesgue measurable sets is a \(\sigma\)-algebra on \(\mathbb{R}\) (2.72). Each interval \((-n,n)\) is open, hence a Borel set, hence Lebesgue measurable. Because a \(\sigma\)-algebra is closed under intersection (2.25), \(A_n = (-n,n) \cap A\) is Lebesgue measurable for every \(n \in \mathbb{Z}^+\). By the forward direction of Exercise 12 applied to \(A_n \subseteq (-n,n)\),

\begin{equation*} |(-n,n) \cap A| + |(-n,n) \setminus A| = |A_n| + |(-n,n) \setminus A_n| = 2n \end{equation*}

for every \(n \in \mathbb{Z}^+\).

Suppose the equality holds for every \(n \in \mathbb{Z}^+\).

Fix \(n \in \mathbb{Z}^+\). As noted above, the hypothesis says \(|A_n| + |(-n,n) \setminus A_n| = n - (-n)\), and \(A_n \subseteq (-n,n)\). By the reverse direction of Exercise 12, \(A_n\) is Lebesgue measurable.

Finally,

\begin{equation*} A = \bigcup_{n=1}^{\infty} A_n , \end{equation*}

because every real number lies in \((-n,n)\) for all sufficiently large \(n\), so every element of \(A\) lies in some \(A_n\); and each \(A_n \subseteq A\). Since \(\mathcal{L}\) is a \(\sigma\)-algebra (2.72) it is closed under countable unions, so \(A\) is Lebesgue measurable.

Problem (2D.14)

Show that \(\frac{1}{4}\) and \(\frac{9}{13}\) are both in the Cantor set.

Solution

By 2.75, the Cantor set \(C\) is the set of numbers in \([0,1]\) that have a base \(3\) representation containing only \(0\)s and \(2\)s. So it suffices to exhibit such a representation for each of the two numbers. We do this by running the usual greedy base \(3\) algorithm (repeatedly multiply by \(3\); the integer part is the next digit, the fractional part is the new remainder) and then verifying the resulting series.

The number \(\frac{1}{4}\).

Running the algorithm:

\begin{equation*} \begin{aligned} 3 \cdot \tfrac{1}{4} &= \tfrac{3}{4} = 0 + \tfrac{3}{4}, &&\text{first digit } 0, \text{ remainder } \tfrac{3}{4};\\ 3 \cdot \tfrac{3}{4} &= \tfrac{9}{4} = 2 + \tfrac{1}{4}, &&\text{second digit } 2, \text{ remainder } \tfrac{1}{4}. \end{aligned} \end{equation*}

The remainder \(\frac{1}{4}\) has reappeared, so the block \(02\) repeats forever, suggesting

\begin{equation*} \tfrac{1}{4} = 0.020202\ldots_3 . \end{equation*}

Verification of this claim directly from the definition of a base \(3\) representation:

\begin{equation*} \sum_{k=1}^{\infty} \frac{2}{3^{2k}} = 2 \sum_{k=1}^{\infty} \Bigl(\frac{1}{9}\Bigr)^{k} = 2 \cdot \frac{\frac{1}{9}}{1 - \frac{1}{9}} = 2 \cdot \frac{1}{8} = \frac{1}{4}. \end{equation*}

Thus \(\frac{1}{4} \in [0,1]\) has a base \(3\) representation whose digits are all \(0\)s and \(2\)s. By 2.75, \(\frac{1}{4} \in C\).

The number \(\frac{9}{13}\).

Running the algorithm:

\begin{equation*} \begin{aligned} 3 \cdot \tfrac{9}{13} &= \tfrac{27}{13} = 2 + \tfrac{1}{13}, &&\text{first digit } 2, \text{ remainder } \tfrac{1}{13};\\ 3 \cdot \tfrac{1}{13} &= \tfrac{3}{13} = 0 + \tfrac{3}{13}, &&\text{second digit } 0, \text{ remainder } \tfrac{3}{13};\\ 3 \cdot \tfrac{3}{13} &= \tfrac{9}{13} = 0 + \tfrac{9}{13}, &&\text{third digit } 0, \text{ remainder } \tfrac{9}{13}. \end{aligned} \end{equation*}

The original remainder \(\frac{9}{13}\) has reappeared, so the block \(200\) repeats forever, suggesting

\begin{equation*} \tfrac{9}{13} = 0.200200200\ldots_3 . \end{equation*}

Verification directly from the definition:

\begin{equation*} \sum_{k=0}^{\infty} \frac{2}{3^{3k+1}} = \frac{2}{3} \sum_{k=0}^{\infty} \Bigl(\frac{1}{27}\Bigr)^{k} = \frac{2}{3} \cdot \frac{1}{1 - \frac{1}{27}} = \frac{2}{3} \cdot \frac{27}{26} = \frac{18}{26} = \frac{9}{13}. \end{equation*}

Thus \(\frac{9}{13} \in [0,1]\) has a base \(3\) representation whose digits are all \(0\)s and \(2\)s. By 2.75, \(\frac{9}{13} \in C\).

Remark. As the book notes just after 2.75, neither of these numbers is an endpoint of any of the intervals removed in the construction of \(C\). Indeed, every such endpoint is of the form \(\frac{m}{3^n}\) for some integers \(m\) and \(n\) (each \(G_n\) is a union of intervals whose endpoints have denominator \(3^n\)), and neither \(\frac{1}{4}\) nor \(\frac{9}{13}\) can be written that way, since in lowest terms their denominators \(4\) and \(13\) are not powers of \(3\). So these are genuine examples of Cantor set elements that the construction does not hand you directly.

Problem (2D.15)

Show that \(\frac{13}{17}\) is not in the Cantor set.

Solution

By 2.75, the Cantor set \(C\) is the set of numbers in \([0,1]\) having a base \(3\) representation whose digits are all \(0\) or \(2\). So it suffices to show that every base \(3\) representation of \(\frac{13}{17}\) contains a digit equal to \(1\).

Step 1: only one base 3 representation.

Suppose \(x \in [0,1]\) has two different base \(3\) representations, say

\begin{equation*} x = \sum_{k=1}^{\infty} \frac{a_k}{3^k} = \sum_{k=1}^{\infty} \frac{b_k}{3^k}, \end{equation*}

with all \(a_k, b_k \in \{0,1,2\}\), and let \(n\) be the smallest index with \(a_n \neq b_n\); say \(a_n < b_n\). Then

\begin{equation*} 0 = \sum_{k=n}^{\infty} \frac{b_k - a_k}{3^k} \geq \frac{b_n - a_n}{3^n} - \sum_{k=n+1}^{\infty} \frac{2}{3^k} = \frac{b_n - a_n}{3^n} - \frac{1}{3^n} \geq 0 , \end{equation*}

because \(b_n - a_n \geq 1\). Hence equality holds throughout, which forces \(b_n - a_n = 1\) and \(b_k - a_k = -2\) for all \(k > n\), i.e. \(b_k = 0\) and \(a_k = 2\) for all \(k > n\). In particular

\begin{equation*} x = \sum_{k=1}^{n} \frac{b_k}{3^k} = \frac{m}{3^n} \quad \text{where } m = \sum_{k=1}^{n} b_k 3^{n-k} \in \mathbb{Z}. \end{equation*}

Thus a number with two base \(3\) representations must be of the form \(m/3^n\).

Now if \(\frac{13}{17} = \frac{m}{3^n}\) for integers \(m, n\), then \(13 \cdot 3^n = 17 m\), so \(17\) divides \(13 \cdot 3^n\). But \(17\) is prime and divides neither \(13\) nor \(3\), a contradiction. Hence \(\frac{13}{17}\) has exactly one base \(3\) representation.

Step 2: that representation contains a 1.

Run the standard digit algorithm: writing \(x_0 = \frac{13}{17}\) and \(a_k = \lfloor 3 x_{k-1} \rfloor\), \(x_k = 3x_{k-1} - a_k\), we get

\begin{equation*} \begin{aligned} 3 \cdot \tfrac{13}{17} &= \tfrac{39}{17} = 2 + \tfrac{5}{17}, &&a_1 = 2,\\ 3 \cdot \tfrac{5}{17} &= \tfrac{15}{17} = 0 + \tfrac{15}{17}, &&a_2 = 0,\\ 3 \cdot \tfrac{15}{17} &= \tfrac{45}{17} = 2 + \tfrac{11}{17}, &&a_3 = 2,\\ 3 \cdot \tfrac{11}{17} &= \tfrac{33}{17} = 1 + \tfrac{16}{17}, &&a_4 = 1 . \end{aligned} \end{equation*}

Each \(x_k\) produced lies in \([0,1)\), so this algorithm does produce a base \(3\) representation of \(\frac{13}{17}\); explicitly

\begin{equation*} \frac{13}{17} = \frac{2}{3} + \frac{0}{3^2} + \frac{2}{3^3} + \frac{1}{3^4} + \frac{1}{3^4}\cdot\frac{16}{17}, \end{equation*}

which one checks directly: \(\frac{61}{81} + \frac{16}{1377} = \frac{1037 + 16}{1377} = \frac{1053}{1377} = \frac{13}{17}\).

So the unique base \(3\) representation of \(\frac{13}{17}\) is \(0.2021\ldots_3\), and its fourth digit is \(1\). By Step 1 there is no other representation, so no base \(3\) representation of \(\frac{13}{17}\) consists only of \(0\)s and \(2\)s. By 2.75, \(\frac{13}{17} \notin C\).

Concretely, \(\frac{13}{17}\) lies in the open interval \(\left(\frac{61}{81}, \frac{62}{81}\right)\), one of the eight intervals making up \(G_4\) (see Exercise 16).

Problem (2D.16)

List the eight open intervals whose union is \(G_4\) in the definition of the Cantor set (2.74).

Solution

By 2.74, \(G_4\) is the union of the middle-third open intervals of the intervals making up \([0,1] \setminus (G_1 \cup G_2 \cup G_3)\). So first we identify those intervals.

From 2.74 we have \(G_1 = \left(\frac{1}{3},\frac{2}{3}\right)\), \(G_2 = \left(\frac{1}{9},\frac{2}{9}\right) \cup \left(\frac{7}{9},\frac{8}{9}\right)\), and

\begin{equation*} G_3 = \left(\tfrac{1}{27},\tfrac{2}{27}\right) \cup \left(\tfrac{7}{27},\tfrac{8}{27}\right) \cup \left(\tfrac{19}{27},\tfrac{20}{27}\right) \cup \left(\tfrac{25}{27},\tfrac{26}{27}\right), \end{equation*}

and \([0,1] \setminus (G_1 \cup G_2) = \left[0,\frac{1}{9}\right] \cup \left[\frac{2}{9},\frac{1}{3}\right] \cup \left[\frac{2}{3},\frac{7}{9}\right] \cup \left[\frac{8}{9},1\right]\). Deleting from each of these four closed intervals of length \(\frac{1}{9}\) its middle third (which is the corresponding interval of \(G_3\)) leaves eight closed intervals of length \(\frac{1}{27}\):

\begin{equation*} \left[0,\tfrac{1}{27}\right],\ \left[\tfrac{2}{27},\tfrac{3}{27}\right],\ \left[\tfrac{6}{27},\tfrac{7}{27}\right],\ \left[\tfrac{8}{27},\tfrac{9}{27}\right],\ \left[\tfrac{18}{27},\tfrac{19}{27}\right],\ \left[\tfrac{20}{27},\tfrac{21}{27}\right],\ \left[\tfrac{24}{27},\tfrac{25}{27}\right],\ \left[\tfrac{26}{27},1\right]. \end{equation*}

These are exactly the intervals of \([0,1] \setminus (G_1 \cup G_2 \cup G_3)\).

The middle third of \(\left[a, a + \frac{1}{27}\right]\) is \(\left(a + \frac{1}{81}, a + \frac{2}{81}\right)\). Writing each left endpoint above as \(\frac{3j}{81}\) with \(j \in \{0,2,6,8,18,20,24,26\}\), the middle third is \(\left(\frac{3j+1}{81}, \frac{3j+2}{81}\right)\). Hence

\begin{equation*} \begin{aligned} G_4 = &\left(\tfrac{1}{81},\tfrac{2}{81}\right) \cup \left(\tfrac{7}{81},\tfrac{8}{81}\right) \cup \left(\tfrac{19}{81},\tfrac{20}{81}\right) \cup \left(\tfrac{25}{81},\tfrac{26}{81}\right)\\ &\cup \left(\tfrac{55}{81},\tfrac{56}{81}\right) \cup \left(\tfrac{61}{81},\tfrac{62}{81}\right) \cup \left(\tfrac{73}{81},\tfrac{74}{81}\right) \cup \left(\tfrac{79}{81},\tfrac{80}{81}\right). \end{aligned} \end{equation*}

As a check, use the base \(3\) description preceding 2.75: there \(G_1 \cup \cdots \cup G_n\) is described as the set of numbers in \([0,1]\) all of whose base \(3\) representations have a \(1\) among the first \(n\) digits, so \(G_4\) consists of those numbers in \([0,1]\) all of whose base \(3\) representations have their first \(1\) in the fourth digit. Such numbers are precisely those in an interval \(\left(\frac{k}{81}, \frac{k+1}{81}\right)\) with \(k = (d_1 d_2 d_3 1)_3 = 27 d_1 + 9 d_2 + 3 d_3 + 1\) and \(d_1, d_2, d_3 \in \{0,2\}\): a point strictly inside such an interval has every base \(3\) representation beginning \(0.d_1 d_2 d_3 1\), since the set of numbers admitting a representation with that prefix is exactly the closed interval \(\left[\frac{k}{81}, \frac{k+1}{81}\right]\). The eight choices of \((d_1,d_2,d_3)\) give

\begin{equation*} k = 1,\ 7,\ 19,\ 25,\ 55,\ 61,\ 73,\ 79, \end{equation*}

in agreement with the list above. Note also that \(|G_4| = 8 \cdot \frac{1}{81} = \frac{8}{81} = \frac{2^{4-1}}{3^4}\), as asserted in the proof of 2.76.

Problem (2D.17)

Let \(C\) denote the Cantor set. Prove that \(\left\{\frac{1}{2}x + \frac{1}{2}y : x, y \in C\right\} = [0,1]\).

Solution

Write \(S = \left\{\frac{1}{2}x + \frac{1}{2}y : x, y \in C\right\}\).

The inclusion \(S \subseteq [0,1]\) is immediate: if \(x, y \in C\) then \(x, y \in [0,1]\) by 2.74, and \([0,1]\) is convex, so \(\frac{1}{2}x + \frac{1}{2}y \in [0,1]\).

For the reverse inclusion, let \(t \in [0,1]\). If \(t = 1\), take \(x = y = 1 \in C\) (indeed \(1 = 0.222\ldots_3\), so \(1 \in C\) by 2.75). So assume \(t \in [0,1)\) and choose a base \(3\) representation

\begin{equation*} t = \sum_{k=1}^{\infty} \frac{t_k}{3^k}, \qquad t_k \in \{0,1,2\}. \end{equation*}

Define digits \(a_k, b_k \in \{0,2\}\) by

\begin{equation*} (a_k, b_k) = \begin{cases} (0,0) & \text{if } t_k = 0,\\ (0,2) & \text{if } t_k = 1,\\ (2,2) & \text{if } t_k = 2. \end{cases} \end{equation*}

In every case \(a_k + b_k = 2 t_k\), and \(a_k, b_k \in \{0,2\}\).

Now set

\begin{equation*} x = \sum_{k=1}^{\infty} \frac{a_k}{3^k}, \qquad y = \sum_{k=1}^{\infty} \frac{b_k}{3^k}. \end{equation*}

Both series converge (they are dominated by \(\sum 2/3^k = 1\)), and both displays exhibit a base \(3\) representation of \(x\) and of \(y\) using only the digits \(0\) and \(2\). Hence \(x \in C\) and \(y \in C\) by 2.75.

Finally, since both series converge absolutely we may add them term by term:

\begin{equation*} \frac{1}{2}x + \frac{1}{2}y = \sum_{k=1}^{\infty} \frac{1}{2} \cdot \frac{a_k + b_k}{3^k} = \sum_{k=1}^{\infty} \frac{t_k}{3^k} = t . \end{equation*}

Thus \(t \in S\). As \(t \in [0,1]\) was arbitrary, \([0,1] \subseteq S\), and therefore \(S = [0,1]\).

This is the counterexample announced in the discussion introducing the Cantor set (just before 2.73): the sum of two sets of Lebesgue measure \(0\) can have positive Lebesgue measure. Indeed, put \(A = \frac{1}{2}C\). Scaling a cover of \(C\) by open intervals scales the sum of the lengths by \(\frac{1}{2}\), so outer measure satisfies \(|A| = \frac{1}{2}|C| = 0\) by 2.76(b); yet \(S = A + A = [0,1]\), so \(|A + A| = 1\).

Problem (2D.18)

Prove that every open interval of \(\mathbb{R}\) contains either infinitely many or no elements in the Cantor set.

Solution

Let \(C\) denote the Cantor set. It suffices to prove that \(C\) has no isolated points, meaning: for every \(x \in C\) and every \(\varepsilon > 0\), the interval \((x - \varepsilon, x + \varepsilon)\) contains infinitely many elements of \(C\). Granting that, let \(I\) be an open interval of \(\mathbb{R}\) that contains at least one element of \(C\), say \(x \in I \cap C\). Because \(I\) is open, there exists \(\varepsilon > 0\) with \((x-\varepsilon, x+\varepsilon) \subseteq I\); then \(I \cap C \supseteq (x-\varepsilon,x+\varepsilon) \cap C\) is infinite. So \(I \cap C\) is either empty or infinite, which is the assertion of the exercise.

Now fix \(x \in C\) and \(\varepsilon > 0\). By 2.75 there is a base \(3\) representation

\begin{equation*} x = \sum_{k=1}^{\infty} \frac{a_k}{3^k}, \qquad a_k \in \{0,2\} \text{ for every } k . \end{equation*}

For each positive integer \(n\) define a new digit string by flipping the \(n\)-th digit:

\begin{equation*} a_k^{(n)} = \begin{cases} a_k & \text{if } k \neq n,\\ 2 - a_n & \text{if } k = n, \end{cases} \qquad\text{and}\qquad x^{(n)} = \sum_{k=1}^{\infty} \frac{a_k^{(n)}}{3^k}. \end{equation*}

Since \(a_n \in \{0,2\}\) we have \(2 - a_n \in \{0,2\}\), so every digit \(a_k^{(n)}\) lies in \(\{0,2\}\). The series converges and its sum lies in \([0,1]\), so the display above is a base \(3\) representation of \(x^{(n)}\) using only \(0\)s and \(2\)s. By 2.75, \(x^{(n)} \in C\).

The two strings differ only in the \(n\)-th digit, so

\begin{equation*} \left|x - x^{(n)}\right| = \frac{|a_n - (2 - a_n)|}{3^n} = \frac{2}{3^n} . \end{equation*}

In particular \(x^{(n)} \neq x\) for every \(n\), and the numbers \(\left|x - x^{(n)}\right| = \frac{2}{3^n}\) are pairwise distinct for distinct \(n\); hence \(x^{(1)}, x^{(2)}, x^{(3)}, \ldots\) are pairwise distinct elements of \(C\), all different from \(x\).

Choose \(N\) with \(\frac{2}{3^N} < \varepsilon\). Then for every \(n \geq N\) we have \(\left|x - x^{(n)}\right| = \frac{2}{3^n} \leq \frac{2}{3^N} < \varepsilon\), so \(x^{(n)} \in (x-\varepsilon, x+\varepsilon)\). Thus \((x - \varepsilon, x + \varepsilon)\) contains the infinitely many distinct points \(x^{(N)}, x^{(N+1)}, x^{(N+2)}, \ldots\) of \(C\), as required.

Combined with 2.76(a), this says that \(C\) is a perfect set: closed with no isolated points. (Consistently with 2.76(c), the conclusion is not that \(C\) contains an interval; the infinitely many nearby Cantor set points found above are separated from \(x\) by gaps belonging to the removed sets \(G_n\).)

Problem (2D.19)

Evaluate \(\int_0^1 \Lambda\), where \(\Lambda\) is the Cantor function.

Solution

The answer is \(\int_0^1 \Lambda = \frac{1}{2}\).

By 2.79 the Cantor function \(\Lambda\) is continuous on \([0,1]\), hence Riemann integrable there; being bounded and continuous it is also Lebesgue measurable, and the two integrals agree, so the computation below settles the value under either reading.

The whole point is the following symmetry.

Claim: \(\Lambda(1-x) = 1 - \Lambda(x)\) for every \(x \in [0,1]\).

Case 1: \(x \in C\). By 2.75 write \(x = \sum_{k=1}^{\infty} a_k 3^{-k}\) with every \(a_k \in \{0,2\}\). This representation is the unique base \(3\) representation of \(x\) using only \(0\)s and \(2\)s: by Step 1 of the solution to Exercise 15, if a number has two base \(3\) representations then in one of them some digit is followed by all \(2\)s while in the other that digit is increased by \(1\) and followed by all \(0\)s; a digit \(d \in \{0,2\}\) increased by \(1\) becomes \(1\) or \(3\), neither of which is an admissible digit in \(\{0,2\}\), so at most one of the two representations consists of \(0\)s and \(2\)s. Since \(\sum_{k=1}^{\infty} 2 \cdot 3^{-k} = 1\),

\begin{equation*} 1 - x = \sum_{k=1}^{\infty} \frac{2 - a_k}{3^k}, \end{equation*}

and every digit \(2 - a_k\) lies in \(\{0,2\}\). So \(1 - x \in C\) and the display is its \(\{0,2\}\)-representation. By 2.77, with \(e_k = a_k/2 \in \{0,1\}\),

\begin{equation*} \Lambda(x) = \sum_{k=1}^{\infty} \frac{e_k}{2^k}, \qquad \Lambda(1-x) = \sum_{k=1}^{\infty} \frac{1 - e_k}{2^k} = \sum_{k=1}^{\infty} \frac{1}{2^k} - \Lambda(x) = 1 - \Lambda(x). \end{equation*}

Case 2: \(x \in [0,1] \setminus C\). By 2.74 we have \([0,1] \setminus C = \bigcup_{n=1}^{\infty} G_n\), so \(x\) lies in some \(G_n\). We first pin down the intervals composing \(G_n\).

Subclaim: for every \(n \geq 1\), the set \([0,1] \setminus (G_1 \cup \cdots \cup G_n)\) is the union of the \(2^n\) closed intervals \(\left[c, c + \frac{1}{3^n}\right]\) with \(c = \sum_{i=1}^{n} \frac{d_i}{3^i}\) and \(d_1, \ldots, d_n \in \{0,2\}\); and \(G_n\) is the union of the \(2^{n-1}\) open intervals \(\left(a, a + \frac{1}{3^n}\right)\) with

\begin{equation*} a = \sum_{i=1}^{n-1} \frac{d_i}{3^i} + \frac{1}{3^n}, \qquad d_1, \ldots, d_{n-1} \in \{0,2\} . \end{equation*}

Proof of the subclaim, by induction on \(n\), where the statement at level \(0\) is read as the assertion that \([0,1]\) is the union of the single closed interval \(\left[0, 1\right]\) attached to the empty digit string (true by inspection). For \(n = 1\): by 2.74, \(G_1 = \left(\frac{1}{3}, \frac{2}{3}\right)\), which is the single interval \(\left(a, a + \frac{1}{3}\right)\) with \(a = \frac{1}{3}\) (an empty sum plus \(\frac{1}{3}\)), and \([0,1] \setminus G_1 = \left[0,\frac{1}{3}\right] \cup \left[\frac{2}{3},1\right]\) is the union of the two intervals attached to \(d_1 = 0\) and \(d_1 = 2\). Now let \(n \geq 2\) and suppose the description of \([0,1] \setminus (G_1 \cup \cdots \cup G_{n-1})\) holds for \(n - 1\), so that set is the union of the intervals \(\left[c, c + \frac{1}{3^{n-1}}\right]\) with \(c = \sum_{i=1}^{n-1} d_i 3^{-i}\), \(d_i \in \{0,2\}\). By 2.74, \(G_n\) is the union of the middle-third open intervals of these, and the middle third of \(\left[c, c + \frac{1}{3^{n-1}}\right]\) is \(\left(c + \frac{1}{3^n},\, c + \frac{2}{3^n}\right)\), which is exactly the interval \(\left(a, a + \frac{1}{3^n}\right)\) with \(a = c + \frac{1}{3^n}\). Removing it leaves the two intervals \(\left[c, c + \frac{1}{3^n}\right]\) and \(\left[c + \frac{2}{3^n}, c + \frac{3}{3^n}\right]\), which are the intervals attached to the digit strings \((d_1, \ldots, d_{n-1}, 0)\) and \((d_1, \ldots, d_{n-1}, 2)\). This proves the subclaim.

So fix the interval \((a,b)\) of \(G_n\) containing \(x\), where \(a\) is as in the subclaim, \(b = a + \frac{1}{3^n}\), and put \(e_i = d_i/2\). Using \(\frac{1}{3^n} = \sum_{k > n} \frac{2}{3^k}\) we get the two \(\{0,2\}\)-representations

\begin{equation*} a = 0.d_1 \ldots d_{n-1} 0\,2\,2\,2\ldots_3, \qquad b = 0.d_1 \ldots d_{n-1} 2\,0\,0\,0\ldots_3 , \end{equation*}

so \(a, b \in C\) and, by 2.77,

\begin{equation*} \Lambda(a) = \sum_{i=1}^{n-1}\frac{e_i}{2^i} + \sum_{k>n}\frac{1}{2^k} = \sum_{i=1}^{n-1}\frac{e_i}{2^i} + \frac{1}{2^n} = \Lambda(b). \end{equation*}

Next, for \(x \in (a,b)\) every base \(3\) representation of \(x\) begins \(0.d_1 \ldots d_{n-1} 1\). Indeed, the set of numbers admitting a representation with a given prefix \(c_1 \ldots c_n\) is exactly the closed interval \(\left[\sum_{i \leq n} c_i 3^{-i},\, \sum_{i \leq n} c_i 3^{-i} + 3^{-n}\right]\), since the tails \(\sum_{k>n} c_k 3^{-k}\) range over all of \(\left[0, 3^{-n}\right]\). Reading a prefix \(c_1 \ldots c_n\) as the base \(3\) integer \(k = \sum_{i \leq n} c_i 3^{n-i} \in \{0, 1, \ldots, 3^n - 1\}\), that interval is \(\left[\frac{k}{3^n}, \frac{k+1}{3^n}\right]\); distinct prefixes give distinct \(k\) and hence intervals with disjoint interiors. The prefix \(d_1 \ldots d_{n-1} 1\) has \(\frac{k}{3^n} = a\), so it gives exactly \([a,b]\), and a point of the open interval \((a,b)\) lies in the interior of \([a,b]\) and therefore in no other prefix interval. Truncating after that first \(1\) and halving the earlier \(2\)s, 2.77 gives \(\Lambda(x) = \sum_{i<n} e_i 2^{-i} + 2^{-n}\), which is the value \(\Lambda(a) = \Lambda(b)\) computed above. Hence \(\Lambda\) is constant on \([a,b]\), with value \(\Lambda(a)\).

Now \(1 - x \in (1-b, 1-a)\), and \(1-a, 1-b \in C\) by Case 1. Since \(\Lambda\) is increasing (2.79),

\begin{equation*} \Lambda(1-b) \leq \Lambda(1-x) \leq \Lambda(1-a), \end{equation*}

while Case 1 gives \(\Lambda(1-a) = 1 - \Lambda(a)\) and \(\Lambda(1-b) = 1 - \Lambda(b) = 1 - \Lambda(a)\). The two outer bounds coincide, so \(\Lambda(1-x) = 1 - \Lambda(a) = 1 - \Lambda(x)\). This proves the claim.

Given the claim, substitute \(u = 1 - x\) in the Riemann integral:

\begin{equation*} \int_0^1 \Lambda(x)\,dx = \int_0^1 \Lambda(1-u)\,du = \int_0^1 \bigl(1 - \Lambda(u)\bigr)\,du = 1 - \int_0^1 \Lambda(u)\,du . \end{equation*}

Writing \(I\) for the common value of these integrals gives \(I = 1 - I\), so \(I = \frac{1}{2}\).

As a check, compute the integral by summing \(\Lambda\) over the intervals it is constant on, whose lengths total \(1\) by the proof of 2.76(b). At stage \(n\) there are \(2^{n-1}\) removed intervals, each of length \(3^{-n}\); by the formula above the value of \(\Lambda\) on the one indexed by \((d_1,\ldots,d_{n-1})\) is \(\frac{2m+1}{2^n}\), where \(m = \sum_{i<n} e_i 2^{n-1-i}\) runs bijectively over \(0, 1, \ldots, 2^{n-1}-1\). Those values sum to \(\sum_{m=0}^{2^{n-1}-1} \frac{2m+1}{2^n} = \frac{(2^{n-1})^2}{2^n} = 2^{n-2}\), so stage \(n\) contributes \(2^{n-2} 3^{-n}\) and the total is

\begin{equation*} \sum_{n=1}^{\infty} \frac{2^{n-2}}{3^n} = \frac{1}{4}\sum_{n=1}^{\infty}\left(\frac{2}{3}\right)^{n} = \frac{1}{4} \cdot 2 = \frac{1}{2}, \end{equation*}

agreeing with the answer above.

Problem (2D.20)

Evaluate each of the following:

(a) \(\Lambda\!\left(\frac{9}{13}\right)\);

(b) \(\Lambda(0.93)\).

Solution

(a) First find the base \(3\) representation of \(\frac{9}{13}\) by the usual digit algorithm:

\begin{equation*} \begin{aligned} 3 \cdot \tfrac{9}{13} &= \tfrac{27}{13} = 2 + \tfrac{1}{13}, &&\text{digit } 2,\\ 3 \cdot \tfrac{1}{13} &= \tfrac{3}{13} = 0 + \tfrac{3}{13}, &&\text{digit } 0,\\ 3 \cdot \tfrac{3}{13} &= \tfrac{9}{13} = 0 + \tfrac{9}{13}, &&\text{digit } 0, \end{aligned} \end{equation*}

and now the remainder \(\frac{9}{13}\) has returned, so the digits repeat with period \(3\):

\begin{equation*} \frac{9}{13} = 0.\overline{200}_3 . \end{equation*}

This is confirmed directly by summing the geometric series: \(\sum_{j=0}^{\infty} \frac{2}{3^{3j+1}} = \frac{2/3}{1 - 1/27} = \frac{2/3 \cdot 27}{26} = \frac{18}{26} = \frac{9}{13}\).

All digits are \(0\) or \(2\), so \(\frac{9}{13} \in C\) by 2.75 (this is the fact asserted in Exercise 14). By 2.77, \(\Lambda\!\left(\frac{9}{13}\right)\) is obtained by replacing each \(2\) by \(1\) and reading the result in base \(2\):

\begin{equation*} \Lambda\!\left(\frac{9}{13}\right) = 0.\overline{100}_2 = \sum_{j=0}^{\infty} \frac{1}{2^{3j+1}} = \frac{1/2}{1 - 1/8} = \frac{4}{7}. \end{equation*}

As a sanity check, \(\frac{2}{3} < \frac{9}{13} < \frac{7}{9}\) (numerically \(0.667 < 0.692 < 0.778\)), and by Example 2.78 \(\Lambda\!\left(\frac{2}{3}\right) = \frac{1}{2}\) while \(\Lambda\) equals \(\frac{3}{4}\) on \(\left(\frac{7}{9},\frac{8}{9}\right)\). Since \(\Lambda\) is increasing (2.79), the value must lie between \(\frac{1}{2}\) and \(\frac{3}{4}\), and indeed \(\frac{4}{7} \approx 0.571\) does.

(b) Here \(0.93\) is a decimal, so we must evaluate \(\Lambda\) at \(\frac{93}{100}\). Its base \(3\) digits:

\begin{equation*} \begin{aligned} 3 \cdot \tfrac{93}{100} &= \tfrac{279}{100} = 2 + \tfrac{79}{100}, &&\text{digit } 2,\\ 3 \cdot \tfrac{79}{100} &= \tfrac{237}{100} = 2 + \tfrac{37}{100}, &&\text{digit } 2,\\ 3 \cdot \tfrac{37}{100} &= \tfrac{111}{100} = 1 + \tfrac{11}{100}, &&\text{digit } 1 . \end{aligned} \end{equation*}

So \(\frac{93}{100} = 0.221\ldots_3\), with a \(1\) appearing in the third digit. Since \(\frac{93}{100}\) is not of the form \(m/3^k\) (that would force \(100 \mid 93 \cdot 3^k\), impossible as \(\gcd(93 \cdot 3^k, 100)\) is \(1\)), by Step 1 of Exercise 15 this base \(3\) representation is unique, and it contains a \(1\); hence \(\frac{93}{100} \notin C\) by 2.75.

By 2.77, truncate after the first \(1\) and replace each \(2\) before that \(1\) by \(1\): the string \(0.221\) becomes \(0.111\), read in base \(2\). Therefore

\begin{equation*} \Lambda(0.93) = 0.111_2 = \frac{1}{2} + \frac{1}{4} + \frac{1}{8} = \frac{7}{8}. \end{equation*}

Equivalently and as a check: \(0.221_3 = \frac{18 + 6 + 1}{27} = \frac{25}{27}\) and \(\frac{26}{27} \approx 0.96296\), so \(\frac{25}{27} < 0.93 < \frac{26}{27}\), i.e. \(0.93\) lies in the interval \(\left(\frac{25}{27},\frac{26}{27}\right)\) of \(G_3\) listed in 2.74. The graph in Example 2.78 shows \(\Lambda\) taking the value \(\frac{7}{8}\) on that rightmost third-stage interval.

Problem (2D.21)

Find each of the following sets:

(a) \(\Lambda^{-1}\bigl(\{\frac{1}{3}\}\bigr)\);

(b) \(\Lambda^{-1}\bigl(\{\frac{5}{16}\}\bigr)\).

Solution

The answers are \(\Lambda^{-1}\bigl(\{\frac{1}{3}\}\bigr) = \left\{\frac{1}{4}\right\}\) and \(\Lambda^{-1}\bigl(\{\frac{5}{16}\}\bigr) = \left[\frac{19}{81}, \frac{20}{81}\right]\).

Set up two bookkeeping facts about 2.77, splitting \([0,1]\) into \(C\) and its complement.

Fact 1 (points of \(C\)). By 2.75 each \(x \in C\) has a base \(3\) representation \(x = \sum_{k} a_k 3^{-k}\) with all \(a_k \in \{0,2\}\), and (as shown in the solution to Exercise 19) that representation is unique. Putting \(e_k = a_k/2 \in \{0,1\}\), definition 2.77 gives \(\Lambda(x) = \sum_{k} e_k 2^{-k}\). Conversely, any string \((e_k)\) of \(0\)s and \(1\)s arises this way from \(x = \sum_k 2 e_k 3^{-k} \in C\). Hence \(x \mapsto (e_k)\) is a bijection between \(C\) and the set of binary digit strings, and under it \(\Lambda(x)\) is the number that the string represents in base \(2\). So for \(y \in [0,1]\), the preimages of \(y\) lying in \(C\) correspond exactly to the base \(2\) representations of \(y\).

Fact 2 (points off \(C\)). By the subclaim in Case 2 of the solution to Exercise 19, together with \([0,1] \setminus C = \bigcup_{n} G_n\) (2.74) and the disjointness of \(G_1, G_2, \ldots\) noted in the proof of 2.76, the set \([0,1] \setminus C\) is the disjoint union, over \(n \geq 1\) and over \((d_1,\ldots,d_{n-1}) \in \{0,2\}^{n-1}\), of the open intervals \((a,b)\) with \(a = \sum_{i<n} d_i 3^{-i} + 3^{-n}\) and \(b = a + 3^{-n}\); and on such an interval \(\Lambda\) has the constant value

\begin{equation*} \sum_{i=1}^{n-1} \frac{e_i}{2^i} + \frac{1}{2^n} = 0.e_1\ldots e_{n-1}1_2 = \frac{2m+1}{2^n}, \qquad e_i = \tfrac{d_i}{2},\ \ m = \sum_{i<n} e_i 2^{n-1-i}. \end{equation*}

As \((d_1,\ldots,d_{n-1})\) ranges over \(\{0,2\}^{n-1}\), the integer \(m\) ranges over \(0,1,\ldots,2^{n-1}-1\) bijectively. So the values \(\Lambda\) takes off \(C\) are exactly the numbers \((2m+1)/2^n\), i.e. the dyadic rationals in \((0,1)\); and such a value determines \(n\) (the \(2\)-adic denominator) and then \(m\), hence determines \((d_1,\ldots,d_{n-1})\) and the interval uniquely.

Fact 3 (base \(2\) representations). The argument of Step 1 in the solution to Exercise 15 works verbatim with \(3\) replaced by \(2\) and the digit bound \(2\) replaced by \(1\): if \(y \in [0,1]\) has two base \(2\) representations, then \(y = m/2^n\) for some integers \(m, n\), and the two representations are \(0.b_1 \ldots b_{n-1} b_n 0 0 0 \ldots_2\) and \(0.b_1 \ldots b_{n-1} (b_n - 1) 1 1 1 \ldots_2\) for a single index \(n\). Conversely, if \(y \in (0,1)\) is dyadic, write \(y = m/2^n\) with \(m\) odd and \(0 < m < 2^n\); then \(y = 0.b_1 \ldots b_n 0 0 0 \ldots_2\) with \(b_n = 1\), and since \(\sum_{k > n} 2^{-k} = 2^{-n}\) we also have \(y = 0.b_1 \ldots b_{n-1} 0 1 1 1 \ldots_2\). So a dyadic rational in \((0,1)\) has exactly two base \(2\) representations, and every other number in \([0,1]\) has exactly one.

(a) \(y = \frac{1}{3}\). If \(\frac{1}{3} = \frac{m}{2^n}\) then \(2^n = 3m\), impossible since \(3 \nmid 2^n\). So \(\frac{1}{3}\) is not a dyadic rational, and by Fact 2 no point outside \(C\) is mapped to it. By Fact 3, \(\frac{1}{3}\) has exactly one base \(2\) representation. That representation is \(0.\overline{01}_2\), since

\begin{equation*} 0.\overline{01}_2 = \sum_{j=1}^{\infty} \frac{1}{4^{j}} = \frac{1/4}{1 - 1/4} = \frac{1}{3}. \end{equation*}

By Fact 1 the one and only preimage is the point of \(C\) whose \(\{0,2\}\)-representation is \(0.\overline{02}_3\), namely

\begin{equation*} x = \sum_{j=1}^{\infty} \frac{2}{9^{j}} = \frac{2/9}{1-1/9} = \frac{2}{8} = \frac{1}{4}. \end{equation*}

Hence \(\Lambda^{-1}\bigl(\{\tfrac{1}{3}\}\bigr) = \left\{\frac{1}{4}\right\}\). (Note \(\frac{1}{4} \in C\), as asserted in Exercise 14.)

(b) \(y = \frac{5}{16}\). This is dyadic, so both cases contribute.

Preimages off \(C\): write \(\frac{5}{16} = \frac{2m+1}{2^n}\); this forces \(n = 4\) and \(2m+1 = 5\), so \(m = 2\), whose \(3\)-digit binary form is \(e_1e_2e_3 = 010\), i.e. \((d_1,d_2,d_3) = (0,2,0)\). By Fact 2 the corresponding interval is \((a,b)\) with

\begin{equation*} a = \frac{0}{3} + \frac{2}{9} + \frac{0}{27} + \frac{1}{81} = \frac{18 + 1}{81} = \frac{19}{81}, \qquad b = a + \frac{1}{81} = \frac{20}{81}, \end{equation*}

which is indeed one of the eight intervals of \(G_4\) found in Exercise 16. So \(\Lambda^{-1}\bigl(\{\tfrac{5}{16}\}\bigr) \setminus C = \left(\frac{19}{81}, \frac{20}{81}\right)\).

Preimages in \(C\): \(\frac{5}{16} = \frac{1}{4} + \frac{1}{16} = 0.0101_2\) is dyadic, so by Fact 3 it has exactly two base \(2\) representations,

\begin{equation*} 0.0101\,000\ldots_2 \qquad\text{and}\qquad 0.0100\,111\ldots_2 . \end{equation*}

By Fact 1 these give the two Cantor set preimages

\begin{equation*} 0.0202\,000\ldots_3 = \frac{2}{9} + \frac{2}{81} = \frac{18+2}{81} = \frac{20}{81}, \end{equation*}

\begin{equation*} 0.0200\,222\ldots_3 = \frac{2}{9} + \sum_{k=5}^{\infty}\frac{2}{3^k} = \frac{18}{81} + \frac{1}{81} = \frac{19}{81}. \end{equation*}

(Both are in \(C\), being the endpoints of the removed interval above, as in Case 2 of Exercise 19.)

Combining,

\begin{equation*} \Lambda^{-1}\left(\left\{\tfrac{5}{16}\right\}\right) = \left\{\tfrac{19}{81}\right\} \cup \left(\tfrac{19}{81}, \tfrac{20}{81}\right) \cup \left\{\tfrac{20}{81}\right\} = \left[\frac{19}{81}, \frac{20}{81}\right]. \end{equation*}

The contrast between (a) and (b) is exactly the dichotomy of Facts 1 and 2: \(\Lambda\) is injective at a value \(y\) that is not a dyadic rational, while at each dyadic rational in \((0,1)\) the level set is a full closed interval, one of the intervals on which the Cantor function is flat.

Problem (2D.22)

(a) Suppose \(x\) is a rational number in \([0, 1]\). Explain why \(\Lambda(x)\) is rational.

(b) Suppose \(x \in C\) is such that \(\Lambda(x)\) is rational. Explain why \(x\) is rational.

Solution

Here \(C\) is the Cantor set and \(\Lambda\) is the Cantor function (2.77). Everything rests on the classical characterization of rationality in terms of digit strings, which we prove first in the exact form needed.

Throughout, fix an integer \(b \ge 2\). By a base \(b\) representation of \(t \in [0,1]\) we mean a sequence \(d_1, d_2, \ldots\) of digits from \(\{0, 1, \ldots, b-1\}\) with

\begin{equation*} t = \sum_{j=1}^{\infty} \frac{d_j}{b^{j}} . \end{equation*}

(The number \(1\) is represented by the constant string \(b-1\); every \(t \in [0,1)\) has one or two representations, as the next lemma shows.)

Lemma 1 (near-uniqueness). If \((c_j)\) and \((c_j^{\prime})\) are base \(b\) representations of the same number \(t\) and \((c_j) \ne (c_j^{\prime})\), then, letting \(n\) be the first index where they differ and assuming without loss of generality that \(c_n > c_n^{\prime}\), we have

\begin{equation*} c_n = c_n^{\prime} + 1, \qquad c_j = 0 \ \text{ and } \ c_j^{\prime} = b - 1 \ \text{ for all } j > n . \end{equation*}

In particular, each of the two strings is eventually constant.

Proof. Since the strings agree before index \(n\),

\begin{equation*} 0 = \sum_{j=1}^{\infty} \frac{c_j - c_j^{\prime}}{b^{j}} = \frac{c_n - c_n^{\prime}}{b^{n}} + \sum_{j > n} \frac{c_j - c_j^{\prime}}{b^{j}} \ge \frac{1}{b^{n}} - \sum_{j > n} \frac{b-1}{b^{j}} = \frac{1}{b^{n}} - \frac{1}{b^{n}} = 0 , \end{equation*}

where the inequality used \(c_n - c_n^{\prime} \ge 1\) and \(c_j - c_j^{\prime} \ge -(b-1)\). Thus every inequality used is an equality, forcing \(c_n - c_n^{\prime} = 1\) and \(c_j - c_j^{\prime} = -(b-1)\), that is \(c_j = 0\) and \(c_j^{\prime} = b-1\), for every \(j > n\). \(\square\)

Lemma 2. Let \((d_j)\) be a base \(b\) representation of \(t \in [0,1]\). Then \(t\) is rational if and only if the sequence \((d_j)\) is eventually periodic (that is, there exist \(N \ge 0\) and \(k \ge 1\) with \(d_{j+k} = d_j\) for all \(j > N\)).

Proof. Suppose first that \((d_j)\) is eventually periodic, with \(d_{j+k} = d_j\) for all \(j > N\). Put

\begin{equation*} z = \sum_{i=1}^{\infty} \frac{d_{N+i}}{b^{i}} , \end{equation*}

so that \(t = \sum_{j \le N} d_j b^{-j} + b^{-N} z\). The string \((d_{N+i})_{i \ge 1}\) is periodic with period \(k\), hence

\begin{equation*} b^{k} z = \sum_{i=1}^{\infty} d_{N+i} b^{\,k-i} = m + \sum_{i=1}^{\infty} \frac{d_{N+k+i}}{b^{i}} = m + z , \end{equation*}

where \(m = \sum_{i=1}^{k} d_{N+i} b^{\,k-i}\) is an integer. Thus \(z = m/(b^{k}-1)\) is rational, and therefore \(t\) is rational.

Conversely, suppose \(t = p/q\) with \(p, q\) integers, \(q \ge 1\). For \(n \ge 0\) let

\begin{equation*} s_n = \sum_{i=1}^{\infty} \frac{d_{n+i}}{b^{i}} = b^{n} t - \sum_{j=1}^{n} d_j b^{\,n-j} . \end{equation*}

Then \(s_n \in [0,1]\), and \(q s_n = b^{n} p - q \sum_{j \le n} d_j b^{\,n-j}\) is an integer. Hence \(q s_n \in \{0, 1, \ldots, q\}\), so the numbers \(s_0, s_1, s_2, \ldots\) take at most \(q+1\) distinct values. By the pigeonhole principle there are \(n < m\) with \(s_n = s_m\). Now \((d_{n+i})_{i \ge 1}\) and \((d_{m+i})_{i \ge 1}\) are two base \(b\) representations of the same number \(s_n = s_m\). If they are equal, then \(d_{j+(m-n)} = d_j\) for every \(j > n\) and \((d_j)\) is eventually periodic. If they are not equal, Lemma 1 says that each of them is eventually constant, so \((d_j)\) is eventually constant and in particular eventually periodic. \(\square\)

(a) Let \(x \in [0,1]\) be rational. There are two cases, according to the two bullet points in the definition 2.77 of \(\Lambda\).

Case 1: \(x \notin C\). By 2.75, a number in \([0,1]\) lies in \(C\) exactly when it has a base \(3\) representation using only the digits \(0\) and \(2\); so every base \(3\) representation of \(x\) contains a digit \(1\). The definition of \(\Lambda\) takes such a representation, truncates it immediately after the first \(1\), replaces each \(2\) occurring before that \(1\) by a \(1\), and reads the resulting finite string \(e_1 e_2 \cdots e_n\) in base \(2\). Hence

\begin{equation*} \Lambda(x) = \sum_{j=1}^{n} \frac{e_j}{2^{j}} , \end{equation*}

a finite sum of rationals, so \(\Lambda(x)\) is rational (indeed a dyadic rational). Note that this case did not use the rationality of \(x\): \(\Lambda\) is rational at every point of \([0,1] \setminus C\), which is unsurprising since \(\Lambda\) is constant on each interval removed in the construction of \(C\).

Case 2: \(x \in C\). Let \(d_1 d_2 \cdots\) be the base \(3\) representation of \(x\) using only \(0\)s and \(2\)s, which is the one used in the definition of \(\Lambda\). Since \(x\) is rational, Lemma 2 (with \(b = 3\)) shows that \((d_j)\) is eventually periodic. By definition,

\begin{equation*} \Lambda(x) = \sum_{j=1}^{\infty} \frac{d_j/2}{2^{j}} , \end{equation*}

that is, \((d_j / 2)_{j \ge 1}\) is a base \(2\) representation of \(\Lambda(x)\). Replacing each digit by half of it does not disturb eventual periodicity, so \((d_j/2)\) is eventually periodic, and Lemma 2 (with \(b = 2\)) shows that \(\Lambda(x)\) is rational.

In both cases \(\Lambda(x)\) is rational, as desired.

(b) Suppose \(x \in C\) and \(\Lambda(x)\) is rational. Let \(d_1 d_2 \cdots\) be the base \(3\) representation of \(x\) consisting only of \(0\)s and \(2\)s; this representation is unique, as asserted in 2.77, because two distinct representations of the same number must, by Lemma 1, differ by exactly \(1\) at the first index where they differ, whereas two digits from \(\{0,2\}\) differ by \(0\) or \(2\).

By the definition of \(\Lambda\), the sequence \((d_j/2)_{j \ge 1}\), whose terms lie in \(\{0,1\}\), is a base \(2\) representation of the number \(\Lambda(x)\). Since \(\Lambda(x)\) is rational, Lemma 2 applied with \(b = 2\) to this particular representation shows that \((d_j/2)\) is eventually periodic. Multiplying every term by \(2\), the sequence \((d_j)\) is eventually periodic. Now Lemma 2 applied with \(b = 3\) shows that

\begin{equation*} x = \sum_{j=1}^{\infty} \frac{d_j}{3^{j}} \end{equation*}

is rational, as desired.

Note that the hypothesis \(x \in C\) cannot be dropped in (b): \(\Lambda\) takes the rational value \(\tfrac12\) at every point of the interval \(\left(\tfrac13, \tfrac23\right)\) by 2.78, and that interval contains irrational numbers.

Problem (2D.23)

Show that there exists a function \(f : \mathbb{R} \to \mathbb{R}\) such that the image under \(f\) of every nonempty open interval is \(\mathbb{R}\).

Solution

We construct such an \(f\) explicitly out of decimal digits. The idea: build a function whose value at \(x\) depends only on the tail of the decimal expansion of \(x\), and which already takes all its values on every set of numbers sharing a fixed initial block of digits. Since every nonempty open interval contains such a set, the function is onto on every interval.

The digit map. For \(x \in \mathbb{R}\) let \(\langle x \rangle = x - \lfloor x \rfloor \in [0,1)\) denote the fractional part. Every \(t \in [0,1)\) has exactly one decimal representation \(t = \sum_{j \ge 1} d_j 10^{-j}\) with digits \(d_j \in \{0, 1, \ldots, 9\}\) that is not eventually equal to \(9\). For existence, set \(a_j = \lfloor 10^{j} t \rfloor\) and \(d_j = a_j - 10 a_{j-1}\) for \(j \ge 1\); from \(a_{j-1} \le 10^{j-1} t < a_{j-1} + 1\) we get \(10 a_{j-1} \le a_j \le 10 a_{j-1} + 9\), so each \(d_j\) is a digit, the sums telescope to \(\sum_{j=1}^{n} d_j 10^{-j} = a_n 10^{-n}\) (using \(a_0 = 0\)), and \(a_n 10^{-n} \le t < a_n 10^{-n} + 10^{-n}\) forces these partial sums to converge to \(t\); moreover \(d_j = 9\) for all \(j > n\) would give \(t = a_n 10^{-n} + 10^{-n}\), contradicting that same inequality. Uniqueness follows from Lemma 1 of Exercise 22 with \(b = 10\): of two distinct representations of one number, one is eventually \(0\) and the other eventually \(9\), so at most one avoids being eventually \(9\). Write \(d_j(x)\) for the \(j\)-th digit of this canonical representation of \(\langle x \rangle\), and define \(\varphi : \mathbb{R} \to [0,9]\) by

\begin{equation*} \varphi(x) = \limsup_{n \to \infty} \frac{d_1(x) + d_2(x) + \cdots + d_n(x)}{n} . \end{equation*}

The averages lie in \([0,9]\), so \(\varphi\) is well defined with values in \([0,9]\).

Composing with a surjection. Define \(h : [0,9] \to \mathbb{R}\) by

\begin{equation*} h(t) = \tan\!\left( \pi \left( \frac{t}{9} - \frac12 \right) \right) \ \text{ for } t \in (0,9), \qquad h(0) = h(9) = 0 . \end{equation*}

As \(t\) runs over \((0,9)\), the quantity \(\pi(t/9 - 1/2)\) runs over \(\left(-\tfrac{\pi}{2}, \tfrac{\pi}{2}\right)\), on which \(\tan\) is a bijection onto \(\mathbb{R}\). Hence \(h\) maps \([0,9]\) onto \(\mathbb{R}\). Set

\begin{equation*} f = h \circ \varphi : \mathbb{R} \to \mathbb{R} . \end{equation*}

We claim \(f(I) = \mathbb{R}\) for every nonempty open interval \(I\). Because \(f(I) \subseteq \mathbb{R}\) and \(h([0,9]) = \mathbb{R}\), it suffices to prove that \(\varphi(I) = [0,9]\).

Step 1: every interval contains a full decimal cylinder. Let \(I\) be a nonempty open interval. Every nonempty open interval contains a nonempty open subinterval lying strictly between two consecutive integers, so we may choose an integer \(m\) and real numbers \(u, v\) with

\begin{equation*} m < u < v < m+1 \qquad \text{and} \qquad (u,v) \subseteq I ; \end{equation*}

then \(\lfloor x \rfloor = m\), and hence \(\langle x \rangle = x - m\), for every \(x \in (u,v)\). Choose \(k \in \mathbb{Z}^{+}\) with \(10^{-k} < (v-u)/2\) and set

\begin{equation*} M = \left\lfloor (u-m) 10^{k} \right\rfloor + 1 , \end{equation*}

an integer satisfying \(M \ge 1\) and \((u-m) 10^{k} < M \le (u-m) 10^{k} + 1\). Dividing by \(10^{k}\) turns these into

\begin{equation*} u - m < \frac{M}{10^{k}} \quad\text{and}\quad \frac{M+1}{10^{k}} \le (u-m) + \frac{2}{10^{k}} < (u-m) + (v-u) = v - m . \end{equation*}

Note the first inequality is strict, which is what will place our points in the open interval. Also \((M+1) 10^{-k} < v - m < 1\), so \(M + 1 < 10^{k}\) and therefore \(1 \le M < 10^{k}\). Let \(c_1 c_2 \cdots c_k\) be the digits of \(M\) written with exactly \(k\) digits, allowing leading zeros, so that \(M 10^{-k} = 0.c_1 c_2 \cdots c_k\). Now let \(e_1, e_2, \ldots\) be any digit sequence that is not eventually equal to \(9\), and put

\begin{equation*} x = m + \frac{M}{10^{k}} + \frac{1}{10^{k}} \sum_{i=1}^{\infty} \frac{e_i}{10^{i}} . \end{equation*}

Because \(\sum_{i \ge 1} e_i 10^{-i} \in [0,1]\), the number \(\langle x \rangle = x - m\) lies in \(\left[ M 10^{-k}, (M+1) 10^{-k} \right]\), and this closed interval is contained in \((u-m, v-m)\) precisely because the two inequalities displayed above are \(u - m < M 10^{-k}\) and \((M+1) 10^{-k} < v-m\). Hence \(x \in (u,v) \subseteq I\). Moreover the canonical decimal digits of \(\langle x \rangle\) are exactly \(c_1, \ldots, c_k, e_1, e_2, \ldots\): this string is a decimal representation of \(\langle x \rangle\), and it is not eventually \(9\) (because \((e_i)\) is not), so it is the canonical one.

A fixed finite prefix does not affect the limit superior of the averages, so

\begin{equation*} \varphi(x) = \limsup_{n \to \infty} \frac{e_1 + \cdots + e_n}{n} . \end{equation*}

Indeed, write \(K = c_1 + \cdots + c_k\) and \(S_j = e_1 + \cdots + e_j\), so that for \(n > k\) the average of the first \(n\) digits of \(\langle x \rangle\) equals \((K + S_{n-k})/n\). Substituting \(j = n-k\) and subtracting,

\begin{equation*} \frac{K + S_j}{j+k} - \frac{S_j}{j} = \frac{K}{j+k} - \frac{k\, S_j}{j (j+k)} , \end{equation*}

and since \(0 \le S_j \le 9j\) the right side has absolute value at most \(\frac{K}{j+k} + \frac{9k}{j+k}\), which tends to \(0\) as \(j \to \infty\). Two sequences whose difference tends to \(0\) have the same limit superior, so \(\varphi(x) = \limsup_{j \to \infty} S_j / j\), as claimed.

Step 2: the tail averages realize every value in \([0,9]\). Fix \(t \in [0,9]\). We produce a digit sequence \((e_i)\), not eventually equal to \(9\), with \(\limsup_n (e_1 + \cdots + e_n)/n = t\).

If \(t = 9\), let \(e_i = 0\) when \(i\) is a power of \(2\) and \(e_i = 9\) otherwise. The number of powers of \(2\) in \(\{1, \ldots, n\}\) is at most \(1 + \log_2 n\), so

\begin{equation*} 9 \ge \frac{e_1 + \cdots + e_n}{n} \ge \frac{9n - 9(1 + \log_2 n)}{n} \longrightarrow 9 , \end{equation*}

whence the limit, and so the limit superior, is \(9\). This sequence has infinitely many \(0\)s, so it is not eventually \(9\).

If \(t \in [0,9)\), define \((e_i)\) greedily: having chosen \(e_1, \ldots, e_{n-1}\), with \(S_{n-1} = e_1 + \cdots + e_{n-1}\) (and \(S_0 = 0\)), put

\begin{equation*} e_n = \begin{cases} 9 & \text{if } \dfrac{S_{n-1} + 9}{n} \le t, \\[4pt] 0 & \text{otherwise.} \end{cases} \end{equation*}

Write \(A_n = S_n / n\). We check \(A_n \le t\) for all \(n \ge 1\) by induction. If \(e_n = 9\) then \(A_n = (S_{n-1}+9)/n \le t\) by the defining rule. If \(e_n = 0\) then \(A_n = S_{n-1}/n \le S_{n-1}/(n-1) = A_{n-1} \le t\) for \(n \ge 2\), while for \(n = 1\) we get \(A_1 = 0 \le t\). So \(A_n \le t\) always, and therefore \(\limsup_n A_n \le t\).

Next, \(e_n = 0\) for infinitely many \(n\): otherwise \(e_n = 9\) for all \(n\) beyond some point, forcing \(A_n \to 9\), contradicting \(A_n \le t < 9\). Whenever \(e_n = 0\) we have \((S_{n-1} + 9)/n > t\), hence

\begin{equation*} A_n = \frac{S_{n-1}}{n} > t - \frac{9}{n} . \end{equation*}

Letting \(n \to \infty\) through the infinitely many indices with \(e_n = 0\) gives \(\limsup_n A_n \ge t\). Therefore \(\limsup_n A_n = t\). Finally, since \(e_n = 0\) infinitely often, \((e_n)\) is not eventually equal to \(9\).

Conclusion. Combining the two steps: given a nonempty open interval \(I\) and any \(t \in [0,9]\), Step 2 supplies a legal digit tail \((e_i)\) with tail average limsup equal to \(t\), and Step 1 converts it into a point \(x \in I\) with \(\varphi(x) = t\). Thus \(\varphi(I) = [0,9]\) and hence

\begin{equation*} f(I) = h(\varphi(I)) = h([0,9]) = \mathbb{R} \end{equation*}

for every nonempty open interval \(I\), as required.

Two remarks. First, such an \(f\) is necessarily discontinuous at every point, since it is unbounded on every interval. Second, and more relevant to this book, the function \(f\) constructed above is nevertheless Borel measurable: each digit function \(d_j\) is constant on each interval of a countable partition of \(\mathbb{R}\) into intervals, hence is Borel measurable; therefore each average \((d_1 + \cdots + d_n)/n\) is Borel measurable, and \(\varphi\), being the limit superior of a sequence of Borel measurable functions, is Borel measurable, since \(\limsup\) is an infimum of suprema and both operations preserve measurability by 2.53. Since \(h\) is Borel measurable (it is continuous on \((0,9)\) and constant on the two-point set \(\{0,9\}\)), the composition \(f = h \circ \varphi\) is Borel measurable by 2.44. So wild behavior of this kind is entirely compatible with measurability.

Problem (2D.24)

For \(A \subseteq \mathbb{R}\), the quantity

\begin{equation*} \sup\{ |F| : F \text{ is a closed bounded subset of } \mathbb{R} \text{ and } F \subseteq A \} \end{equation*}

is called the inner measure of \(A\).

(a) Show that if \(A\) is a Lebesgue measurable subset of \(\mathbb{R}\), then the inner measure of \(A\) equals the outer measure of \(A\).

(b) Show that inner measure is not a measure on the \(\sigma\)-algebra of all subsets of \(\mathbb{R}\).

Solution

Write \(\mu_{*}(A)\) for the inner measure of \(A\) and, as in the book, \(|A|\) for the outer measure of \(A\). Note that the empty set is a closed bounded subset of every \(A\), so the supremum is over a nonempty collection and \(\mu_{*}(A) \in [0, \infty]\) is always defined. Note also that every closed subset of \(\mathbb{R}\) is a Borel set and hence Lebesgue measurable, so \(|F|\) in the definition is the Lebesgue measure of \(F\).

(a) Let \(A \subseteq \mathbb{R}\) be Lebesgue measurable.

The inequality \(\mu_{*}(A) \le |A|\). This holds for every set \(A\), measurable or not: if \(F \subseteq A\) then \(|F| \le |A|\) because outer measure preserves order (2.5). Taking the supremum over all closed bounded \(F \subseteq A\) gives \(\mu_{*}(A) \le |A|\).

The inequality \(\mu_{*}(A) \ge |A|\). Let \(\varepsilon > 0\). Because \(A\) is Lebesgue measurable, 2.71 (the implication (a) \(\Rightarrow\) (b)) provides a closed set \(F \subseteq A\) with

\begin{equation*} |A \setminus F| < \varepsilon . \end{equation*}

The set \(F\) need not be bounded, so we truncate it. For \(n \in \mathbb{Z}^{+}\) put

\begin{equation*} F_n = F \cap [-n, n] . \end{equation*}

Each \(F_n\) is closed (an intersection of two closed sets) and bounded, and \(F_n \subseteq F \subseteq A\); hence \(|F_n| \le \mu_{*}(A)\) for every \(n\). The sets \(F_1 \subseteq F_2 \subseteq \cdots\) are Lebesgue measurable with union \(F\), so by 2.59 (the measure of an increasing union, applied to Lebesgue measure, which is a measure on the \(\sigma\)-algebra of Lebesgue measurable sets by 2.72)

\begin{equation*} \lim_{n \to \infty} |F_n| = |F| . \end{equation*}

Consequently \(\sup_n |F_n| = |F|\) and therefore

\begin{equation*} \mu_{*}(A) \ge |F| . \end{equation*}

Now \(A = F \cup (A \setminus F)\), so subadditivity of outer measure (2.8) gives

\begin{equation*} |A| \le |F| + |A \setminus F| < |F| + \varepsilon . \end{equation*}

If \(|A| = \infty\), this forces \(|F| = \infty\) and hence \(\mu_{*}(A) = \infty = |A|\), and we are done. If \(|A| < \infty\), we obtain

\begin{equation*} \mu_{*}(A) \ge |F| > |A| - \varepsilon . \end{equation*}

Since \(\varepsilon > 0\) was arbitrary, \(\mu_{*}(A) \ge |A|\).

Combining the two inequalities, \(\mu_{*}(A) = |A|\) whenever \(A\) is Lebesgue measurable. In particular \(\mu_{*}([0,1)) = |[0,1)| = 1\), a fact used below.

(b) Inner measure does satisfy \(\mu_{*}(\varnothing) = 0\), so to show that it is not a measure on the \(\sigma\)-algebra of all subsets of \(\mathbb{R}\) we must exhibit a sequence of disjoint sets whose inner measures do not add up to the inner measure of their union. We use the Vitali-type construction that appears in the discussion following 2.72.

Two preliminary observations. First, \(\mu_{*}\) is translation invariant: if \(t \in \mathbb{R}\) and \(A \subseteq \mathbb{R}\), then \(F \mapsto t + F\) is a bijection from the closed bounded subsets of \(A\) onto the closed bounded subsets of \(t + A\), and \(|t + F| = |F|\) by 2.7; hence \(\mu_{*}(t + A) = \mu_{*}(A)\). Second, \(\mu_{*}\) is monotone: \(A \subseteq B\) implies \(\mu_{*}(A) \le \mu_{*}(B)\), since every closed bounded subset of \(A\) is one of \(B\).

Rotation of \([0,1)\). For \(r \in [0,1)\) and \(S \subseteq [0,1)\) define

\begin{equation*} S \oplus r = \bigl( (S \cap [0, 1-r)) + r \bigr) \cup \bigl( (S \cap [1-r, 1)) + r - 1 \bigr) \subseteq [0,1) , \end{equation*}

so \(x \mapsto x \oplus r\) (addition modulo \(1\)) is a bijection of \([0,1)\) onto itself. Note that if \(S\) is Lebesgue measurable then so is \(S \oplus r\), and, since the two pieces above are disjoint, additivity of Lebesgue measure (2.72) together with translation invariance of outer measure (2.7) gives

\begin{equation*} |S \oplus r| = |S \cap [0,1-r)| + |S \cap [1-r,1)| = |S| . \end{equation*}

The Vitali set. Define an equivalence relation on \([0,1)\) by \(x \sim y\) if and only if \(x - y \in \mathbb{Q}\). By the axiom of choice, let \(V \subseteq [0,1)\) contain exactly one element of each equivalence class. Let \(r_1, r_2, r_3, \ldots\) be an enumeration of the countably infinite set \(\mathbb{Q} \cap [0,1)\), and set

\begin{equation*} V_k = V \oplus r_k \qquad (k \in \mathbb{Z}^{+}) . \end{equation*}

The sets \(V_1, V_2, \ldots\) are disjoint with union \([0,1)\). Suppose \(x \in V_j \cap V_k\). Then \(x = v \oplus r_j = w \oplus r_k\) for some \(v, w \in V\); since \(x - (v + r_j) \in \{0, -1\}\) and likewise \(x - (w + r_k) \in \{0, -1\}\), we get \(v - w \in \mathbb{Q}\), so \(v \sim w\) and therefore \(v = w\) by the choice of \(V\). Then \(v \oplus r_j = v \oplus r_k\), and since \(r \mapsto v \oplus r\) is injective on \([0,1)\) we get \(r_j = r_k\), i.e. \(j = k\). For the union, let \(x \in [0,1)\) and let \(v \in V\) be the representative of the class of \(x\). Then \(x - v \in \mathbb{Q} \cap (-1, 1)\); put \(r = x - v\) if \(x \ge v\) and \(r = x - v + 1\) otherwise, so \(r \in \mathbb{Q} \cap [0,1)\) and \(x = v \oplus r\). Thus \(x \in V_k\) for the index \(k\) with \(r_k = r\).

Each \(V_k\) has inner measure \(0\). Since rotation by \(r_k\) maps each equivalence class onto itself (because \(x \oplus r_k \sim x\) for rational \(r_k\)) and is injective, the set \(V_k\), like \(V\), contains exactly one element of each equivalence class. So it suffices to prove the following: if \(S \subseteq [0,1)\) contains at most one element of each equivalence class, then \(\mu_{*}(S) = 0\).

Let \(F \subseteq S\) be closed and bounded. Then \(F\) is a Borel set, hence Lebesgue measurable. The sets \(F \oplus r\), for \(r \in \mathbb{Q} \cap [0,1)\), are pairwise disjoint: if \(x \oplus r = y \oplus r^{\prime}\) with \(x, y \in F\) and \(r \ne r^{\prime}\), then \(x - y \in \mathbb{Q}\), so \(x \sim y\); also \(x \ne y\), since \(x = y\) would give \(x \oplus r = x \oplus r^{\prime}\) and hence \(r = r^{\prime}\); but \(x, y \in F \subseteq S\) and \(S\) contains at most one element of each class, a contradiction. Each \(F \oplus r\) is a Lebesgue measurable subset of \([0,1)\) with \(|F \oplus r| = |F|\), by the computation above. Therefore, by countable additivity of Lebesgue measure (2.72) and monotonicity (2.5),

\begin{equation*} \sum_{r \in \mathbb{Q} \cap [0,1)} |F| = \sum_{r \in \mathbb{Q} \cap [0,1)} |F \oplus r| = \Bigl| \bigcup_{r \in \mathbb{Q} \cap [0,1)} (F \oplus r) \Bigr| \le |[0,1)| = 1 . \end{equation*}

The left side is a sum of countably infinitely many copies of the same number \(|F|\), so it equals \(\infty\) unless \(|F| = 0\). Hence \(|F| = 0\). As \(F\) was an arbitrary closed bounded subset of \(S\), we conclude \(\mu_{*}(S) = 0\), and in particular

\begin{equation*} \mu_{*}(V_k) = 0 \qquad \text{for every } k \in \mathbb{Z}^{+} . \end{equation*}

Conclusion. The sets \(V_1, V_2, \ldots\) are disjoint subsets of \(\mathbb{R}\), all belonging to the \(\sigma\)-algebra of all subsets of \(\mathbb{R}\), and their union is \([0,1)\). By part (a),

\begin{equation*} \mu_{*}\Bigl( \bigcup_{k=1}^{\infty} V_k \Bigr) = \mu_{*}\bigl( [0,1) \bigr) = \bigl| [0,1) \bigr| = 1 , \end{equation*}

whereas

\begin{equation*} \sum_{k=1}^{\infty} \mu_{*}(V_k) = \sum_{k=1}^{\infty} 0 = 0 . \end{equation*}

Thus \(\mu_{*}\) fails to be countably additive, so inner measure is not a measure on the \(\sigma\)-algebra of all subsets of \(\mathbb{R}\).

Contrast this with part (a): restricted to the \(\sigma\)-algebra of Lebesgue measurable sets, \(\mu_{*}\) coincides with Lebesgue measure and hence is a measure there. The failure in (b) is caused precisely by the nonmeasurable sets, which can have positive outer measure while containing no closed bounded set of positive measure.

Exercises 2E

Problem (2E.1)

Suppose \(X\) is a finite set. Explain why a sequence of functions from \(X\) to \(\mathbb{R}\) that converges pointwise on \(X\) also converges uniformly on \(X\).

Solution

If \(X = \emptyset\) the assertion is vacuous (every condition of the form “for all \(x \in X\)” holds), so assume \(X = \{x_1, \dots, x_N\}\) with \(N \in \mathbb{Z}^+\).

Suppose \(f_1, f_2, \dots\) is a sequence of functions from \(X\) to \(\mathbb{R}\) converging pointwise on \(X\) to \(f : X \to \mathbb{R}\). Let \(\varepsilon > 0\). By the definition of pointwise convergence (2.82), for each \(j \in \{1, \dots, N\}\) there exists \(n_j \in \mathbb{Z}^+\) such that

\begin{equation*} |f_k(x_j) - f(x_j)| < \varepsilon \quad \text{for all integers } k \ge n_j . \end{equation*}

Because \(\{n_1, \dots, n_N\}\) is a finite set of positive integers, it has a maximum. Let

\begin{equation*} n = \max\{n_1, \dots, n_N\} \in \mathbb{Z}^+ . \end{equation*}

Now suppose \(k \ge n\) is an integer and \(x \in X\). Then \(x = x_j\) for some \(j \in \{1, \dots, N\}\), and \(k \ge n \ge n_j\), so \(|f_k(x) - f(x)| = |f_k(x_j) - f(x_j)| < \varepsilon\).

Thus for every \(\varepsilon > 0\) we have produced a single \(n \in \mathbb{Z}^+\) that works simultaneously for all \(x \in X\). By 2.82, \(f_1, f_2, \dots\) converges uniformly on \(X\) to \(f\).

The point is that pointwise convergence differs from uniform convergence only in the order of the quantifiers: pointwise convergence supplies a threshold \(n\) depending on both \(\varepsilon\) and \(x\), while uniform convergence demands a threshold depending only on \(\varepsilon\). When \(X\) is finite, the supremum over \(x \in X\) of these thresholds is attained and is itself a positive integer, so the two notions coincide. (The argument fails for infinite \(X\) because a supremum of infinitely many positive integers may be \(\infty\); see Exercise 2 in this section.)

Problem (2E.2)

Give an example of a sequence of functions from \(\mathbb{Z}^+\) to \(\mathbb{R}\) that converges pointwise on \(\mathbb{Z}^+\) but does not converge uniformly on \(\mathbb{Z}^+\).

Solution

For each \(k \in \mathbb{Z}^+\) define \(f_k : \mathbb{Z}^+ \to \mathbb{R}\) by

\begin{equation*} f_k(m) = \begin{cases} 1 & \text{if } m = k, \\ 0 & \text{if } m \neq k, \end{cases} \end{equation*}

in other words \(f_k = \chi_{\{k\}}\).

Pointwise convergence: fix \(m \in \mathbb{Z}^+\). For every integer \(k > m\) we have \(k \neq m\) and hence \(f_k(m) = 0\). Thus \(\lim_{k \to \infty} f_k(m) = 0\). Since \(m\) was arbitrary, \(f_1, f_2, \dots\) converges pointwise on \(\mathbb{Z}^+\) to the function \(f\) that is identically \(0\).

Failure of uniform convergence: if \(f_1, f_2, \dots\) converged uniformly on \(\mathbb{Z}^+\), the uniform limit would have to agree with the pointwise limit \(f = 0\) (uniform convergence implies pointwise convergence). Take \(\varepsilon = \tfrac12\). For every \(n \in \mathbb{Z}^+\) the integer \(k = n\) satisfies \(k \ge n\), and the point \(x = n \in \mathbb{Z}^+\) satisfies

\begin{equation*} |f_k(x) - f(x)| = |f_n(n) - 0| = 1 > \tfrac12 . \end{equation*}

So no \(n \in \mathbb{Z}^+\) has the property required in 2.82 for \(\varepsilon = \tfrac12\). Hence \(f_1, f_2, \dots\) does not converge uniformly on \(\mathbb{Z}^+\).

Equivalently, \(\sup_{m \in \mathbb{Z}^+} |f_k(m) - f(m)| = 1\) for every \(k\), so the sequence cannot converge uniformly to \(f\).

Note how this example is exactly the failure of the argument in Exercise 1: for the point \(m\), the threshold needed is \(n_m = m + 1\), and \(\sup_m n_m = \infty\).

Problem (2E.3)

Give an example of a sequence of continuous functions \(f_1, f_2, \dots\) from \([0,1]\) to \(\mathbb{R}\) that converges pointwise to a function \(f : [0,1] \to \mathbb{R}\) that is not a bounded function.

Solution

For each \(k \in \mathbb{Z}^+\) define \(f_k : [0,1] \to \mathbb{R}\) by

\begin{equation*} f_k(x) = \begin{cases} k^2 x & \text{if } 0 \le x \le \frac{1}{k}, \\ \dfrac{1}{x} & \text{if } \frac{1}{k} < x \le 1. \end{cases} \end{equation*}

Thus the graph of \(f_k\) rises linearly from \((0,0)\) to \(\bigl(\frac1k, k\bigr)\) and then follows the hyperbola \(y = 1/x\) out to \((1,1)\).

Each \(f_k\) is continuous: it is continuous on \([0, \frac1k]\) and on \([\frac1k, 1]\) (on the latter it is \(1/x\), continuous since \(x \ge \frac1k > 0\)), and the two formulas agree at the junction, since \(k^2 \cdot \frac1k = k = \frac{1}{1/k}\). A function that is continuous on each of two closed intervals whose union is \([0,1]\) and which agree at the shared endpoint is continuous on \([0,1]\).

Define \(f : [0,1] \to \mathbb{R}\) by

\begin{equation*} f(x) = \begin{cases} 0 & \text{if } x = 0, \\ \dfrac{1}{x} & \text{if } 0 < x \le 1. \end{cases} \end{equation*}

Note that \(f\) takes values in \(\mathbb{R}\) (it is real-valued at every point of \([0,1]\), including \(0\)).

The sequence converges pointwise to \(f\): if \(x = 0\), then \(f_k(0) = k^2 \cdot 0 = 0 = f(0)\) for every \(k\), so \(\lim_{k \to \infty} f_k(0) = f(0)\). If \(0 < x \le 1\), choose \(n \in \mathbb{Z}^+\) with \(n > 1/x\); then for every integer \(k \ge n\) we have \(\frac1k \le \frac1n < x\), hence \(f_k(x) = \frac1x = f(x)\). Thus \(\lim_{k \to \infty} f_k(x) = f(x)\) (the sequence is eventually constant at each such \(x\)).

Finally, \(f\) is not bounded: for each \(k \in \mathbb{Z}^+\) the point \(x = \frac1k \in [0,1]\) satisfies \(f(x) = k\), so \(\sup_{x \in [0,1]} |f(x)| = \infty\).

This also gives another example for Exercise 2’s phenomenon: the convergence here is certainly not uniform, since a uniform limit of bounded functions is bounded, while each \(f_k\) is bounded (by \(k\)) and \(f\) is not.

Problem (2E.4)

Prove or give a counterexample: If \(A \subseteq \mathbb{R}\) and \(f_1, f_2, \dots\) is a sequence of uniformly continuous functions from \(A\) to \(\mathbb{R}\) that converges uniformly to a function \(f : A \to \mathbb{R}\), then \(f\) is uniformly continuous on \(A\).

Solution

The statement is true. This is the uniform-continuity analogue of 2.84 (which says that a uniform limit of functions continuous at a point is continuous at that point), and it is proved by the same \(\varepsilon/3\) device, with the key observation that the \(\delta\) produced for \(f_n\) does not depend on the point.

Recall that \(g : A \to \mathbb{R}\) is uniformly continuous on \(A\) if for every \(\varepsilon > 0\) there exists \(\delta > 0\) such that \(|g(x) - g(y)| < \varepsilon\) for all \(x, y \in A\) with \(|x - y| < \delta\).

Suppose \(\varepsilon > 0\). Because \(f_1, f_2, \dots\) converges uniformly on \(A\) to \(f\), there exists \(n \in \mathbb{Z}^+\) such that

\begin{equation*} |f_n(z) - f(z)| < \frac{\varepsilon}{3} \quad \text{for all } z \in A . \end{equation*}

Because \(f_n\) is uniformly continuous on \(A\), there exists \(\delta > 0\) such that

\begin{equation*} |f_n(x) - f_n(y)| < \frac{\varepsilon}{3} \quad \text{for all } x, y \in A \text{ with } |x - y| < \delta . \end{equation*}

Now suppose \(x, y \in A\) and \(|x - y| < \delta\). Then

\begin{equation*} \begin{aligned} |f(x) - f(y)| &\le |f(x) - f_n(x)| + |f_n(x) - f_n(y)| + |f_n(y) - f(y)| \\ &< \frac{\varepsilon}{3} + \frac{\varepsilon}{3} + \frac{\varepsilon}{3} \\ &= \varepsilon . \end{aligned} \end{equation*}

Thus for every \(\varepsilon > 0\) we have produced \(\delta > 0\) that works for all pairs of points of \(A\) simultaneously. Hence \(f\) is uniformly continuous on \(A\).

Two remarks on where the hypotheses are used. Uniform convergence is essential: the sequence of Exercise 3 above consists of functions each of which is uniformly continuous on \([0,1]\) (each \(f_k\) is continuous on the compact interval \([0,1]\), hence uniformly continuous there), yet the pointwise limit \(f\) is not even bounded, so certainly not uniformly continuous. And note that nothing about \(A\) was needed: the argument works verbatim for functions on any metric space, since the only property of \(A\) used is that \(|x-y|\) makes sense.

Problem (2E.5)

Give an example to show that Egorov’s Theorem can fail without the hypothesis that \(\mu(X) < \infty\).

Solution

Let \(X = \mathbb{R}\), let \(\mathcal{S}\) be the \(\sigma\)-algebra of Borel subsets of \(\mathbb{R}\), and let \(\mu\) be Lebesgue measure on \(\mathcal{S}\) (so \(\mu(E) = |E|\), the outer measure of \(E\)). Then \(\mu(X) = \infty\).

For each \(k \in \mathbb{Z}^+\) define \(f_k : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} f_k = \chi_{[k,\, k+1]} . \end{equation*}

Each \(f_k\) is \(\mathcal{S}\)-measurable, since \([k, k+1]\) is a Borel set (the characteristic function of a set \(E\) is \(\mathcal{S}\)-measurable if and only if \(E \in \mathcal{S}\)).

The sequence converges pointwise on \(\mathbb{R}\) to the function \(f\) that is identically \(0\): given \(x \in \mathbb{R}\), every integer \(k > x\) satisfies \(x \notin [k, k+1]\), so \(f_k(x) = 0\) for all such \(k\).

Now take \(\varepsilon = 1\) and suppose \(E \in \mathcal{S}\) satisfies \(\mu(\mathbb{R} \setminus E) < 1\). Fix any \(k \in \mathbb{Z}^+\). Since \([k, k+1] \setminus E \subseteq \mathbb{R} \setminus E\), monotonicity of outer measure gives

\begin{equation*} \bigl| [k, k+1] \setminus E \bigr| \le \bigl| \mathbb{R} \setminus E \bigr| < 1 = \bigl| [k, k+1] \bigr| . \end{equation*}

Hence \([k, k+1] \setminus E \neq [k, k+1]\), which means \([k, k+1] \cap E \neq \emptyset\). Choose \(x_k \in [k, k+1] \cap E\). Then

\begin{equation*} |f_k(x_k) - f(x_k)| = 1 . \end{equation*}

Consequently, for the value \(\varepsilon^{\prime} = \tfrac12\) in the definition of uniform convergence there is no \(n \in \mathbb{Z}^+\) that works: whatever \(n\) is chosen, the integer \(k = n\) and the point \(x_n \in E\) give \(|f_n(x_n) - f(x_n)| = 1 > \tfrac12\). So \(f_1, f_2, \dots\) does not converge uniformly to \(f\) on \(E\).

Since \(E\) was an arbitrary member of \(\mathcal{S}\) with \(\mu(\mathbb{R} \setminus E) < 1\), there is no set \(E \in \mathcal{S}\) with \(\mu(X \setminus E) < \varepsilon\) on which the convergence is uniform. Thus the conclusion of Egorov’s Theorem (2.85) fails for \(\varepsilon = 1\), even though every other hypothesis of 2.85 holds. The hypothesis \(\mu(X) < \infty\) therefore cannot be dropped.

(It is instructive to see exactly where the proof of 2.85 breaks. There 2.59 supplies \(\lim_{m \to \infty} \mu(A_{m,n}) = \mu(X)\), and finiteness of \(\mu(X)\) is what allows this to be turned into the choice of \(m_n\) with \(\mu(X) - \mu(A_{m_n,n}) < \varepsilon/2^n\), and then into \(\mu(X \setminus A_{m_n,n}) < \varepsilon/2^n\) by 2.57(b). Both steps are meaningless when \(\mu(X) = \infty\). Indeed, in the example above

\begin{equation*} A_{m,n} = \bigcap_{k \ge m} \bigl\{ x \in \mathbb{R} : |f_k(x)| < \tfrac1n \bigr\} = \mathbb{R} \setminus \bigcup_{k \ge m} [k, k+1] = (-\infty, m), \end{equation*}

so \(\mu(A_{m,n}) = \infty = \mu(X)\) for every \(m\) and \(n\), while \(\mu(\mathbb{R} \setminus A_{m,n}) = \mu([m,\infty)) = \infty\) for every \(m\) and \(n\). The sets that the proof would select are never of small complementary measure.)

Problem (2E.6)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space with \(\mu(X) < \infty\). Suppose \(f_1, f_2, \dots\) is a sequence of \(\mathcal{S}\)-measurable functions from \(X\) to \(\mathbb{R}\) such that \(\lim_{k \to \infty} f_k(x) = \infty\) for each \(x \in X\). Prove that for every \(\varepsilon > 0\), there exists a set \(E \in \mathcal{S}\) such that \(\mu(X \setminus E) < \varepsilon\) and \(f_1, f_2, \dots\) converges uniformly to \(\infty\) on \(E\) (meaning that for every \(t > 0\), there exists \(n \in \mathbb{Z}^+\) such that \(f_k(x) > t\) for all integers \(k \ge n\) and all \(x \in E\)).

[The exercise above is an Egorov-type theorem for sequences of functions that converge pointwise to \(\infty\).]

Solution

The proof follows the proof of Egorov’s Theorem (2.85), with the sets \(\{x : |f_k(x) - f(x)| < \frac1n\}\) replaced by the sets \(\{x : f_k(x) > n\}\).

Suppose \(\varepsilon > 0\). Temporarily fix \(n \in \mathbb{Z}^+\). For \(m \in \mathbb{Z}^+\) let

\begin{equation*} A_{m,n} = \bigcap_{k = m}^{\infty} \{ x \in X : f_k(x) > n \} . \end{equation*}

Each set \(\{x \in X : f_k(x) > n\}\) equals \(f_k^{-1}\bigl((n, \infty)\bigr)\), which belongs to \(\mathcal{S}\) because \(f_k\) is \(\mathcal{S}\)-measurable and \((n, \infty)\) is a Borel subset of \(\mathbb{R}\). Hence each \(A_{m,n} \in \mathcal{S}\), being a countable intersection of elements of the \(\sigma\)-algebra \(\mathcal{S}\).

Clearly \(A_{1,n} \subseteq A_{2,n} \subseteq \cdots\), because intersecting over fewer indices gives a larger set. Moreover

\begin{equation*} \bigcup_{m = 1}^{\infty} A_{m,n} = X . \end{equation*}

Indeed, the inclusion \(\subseteq\) is trivial; for the reverse, let \(x \in X\). Since \(\lim_{k \to \infty} f_k(x) = \infty\), there exists \(m \in \mathbb{Z}^+\) such that \(f_k(x) > n\) for all integers \(k \ge m\), which says exactly that \(x \in A_{m,n}\).

By 2.59 (measure of an increasing union), \(\lim_{m \to \infty} \mu(A_{m,n}) = \mu(X)\). Since \(\mu(X) < \infty\), there exists \(m_n \in \mathbb{Z}^+\) such that

\begin{equation*} \mu(X) - \mu(A_{m_n, n}) < \frac{\varepsilon}{2^n} . \end{equation*}

Because \(A_{m_n,n} \subseteq X\) and \(\mu(X) < \infty\), 2.57(b) gives \(\mu(X \setminus A_{m_n,n}) = \mu(X) - \mu(A_{m_n,n}) < \varepsilon/2^n\) (the subtraction is legitimate precisely because \(\mu(A_{m_n,n}) \le \mu(X) < \infty\)).

Now let

\begin{equation*} E = \bigcap_{n = 1}^{\infty} A_{m_n, n} . \end{equation*}

Then \(E \in \mathcal{S}\), and by countable subadditivity (2.58),

\begin{equation*} \mu(X \setminus E) = \mu\Bigl( \bigcup_{n=1}^{\infty} (X \setminus A_{m_n, n}) \Bigr) \le \sum_{n=1}^{\infty} \mu(X \setminus A_{m_n, n}) < \sum_{n=1}^{\infty} \frac{\varepsilon}{2^n} = \varepsilon . \end{equation*}

It remains to verify that \(f_1, f_2, \dots\) converges uniformly to \(\infty\) on \(E\). Suppose \(t > 0\). Choose \(n \in \mathbb{Z}^+\) with \(n \ge t\), and set \(n_0 = m_n\). Because \(E \subseteq A_{m_n, n}\), every \(x \in E\) lies in \(A_{m_n,n}\), which means

\begin{equation*} f_k(x) > n \ge t \quad \text{for all integers } k \ge m_n = n_0 \text{ and all } x \in E . \end{equation*}

Thus for every \(t > 0\) there exists \(n_0 \in \mathbb{Z}^+\) with \(f_k(x) > t\) for all integers \(k \ge n_0\) and all \(x \in E\), which is exactly the assertion that \(f_1, f_2, \dots\) converges uniformly to \(\infty\) on \(E\).

As with Egorov’s Theorem itself, the hypothesis \(\mu(X) < \infty\) cannot be omitted. Let \(X = \mathbb{R}\) with \(\mathcal{S}\) the Borel sets and \(\mu\) Lebesgue measure, and for \(k \in \mathbb{Z}^+\) define the continuous (hence Borel measurable) function

\begin{equation*} f_k(x) = \max\{ k - |x|,\, 0 \} . \end{equation*}

For each fixed \(x \in \mathbb{R}\) we have \(f_k(x) = k - |x|\) once \(k > |x|\), so \(\lim_{k \to \infty} f_k(x) = \infty\). Now take \(\varepsilon = 1\) and let \(E \in \mathcal{S}\) satisfy \(\mu(\mathbb{R} \setminus E) < 1\). Exactly as in Exercise 5, for each \(k \in \mathbb{Z}^+\) the inequality \(\bigl| [k, k+1] \setminus E \bigr| \le \bigl| \mathbb{R} \setminus E \bigr| < 1 = \bigl| [k, k+1] \bigr|\) forces \(E \cap [k, k+1] \neq \emptyset\); choose \(x_k\) in that intersection. Since \(|x_k| \ge k\), we get \(f_k(x_k) = 0\). Hence for \(t = 1\) there is no \(n \in \mathbb{Z}^+\) with \(f_k(x) > 1\) for all \(k \ge n\) and all \(x \in E\), because \(k = n\) and \(x = x_n \in E\) give \(f_n(x_n) = 0\). So \(f_1, f_2, \dots\) converges uniformly to \(\infty\) on no such \(E\), and the conclusion fails for \(\varepsilon = 1\).

Problem (2E.7)

Suppose \(F\) is a closed bounded subset of \(\mathbb{R}\) and \(g_1, g_2, \dots\) is an increasing sequence of continuous real-valued functions on \(F\) (thus \(g_1(x) \le g_2(x) \le \cdots\) for all \(x \in F\)) such that \(\sup\{g_1(x), g_2(x), \dots\} < \infty\) for each \(x \in F\). Define a real-valued function \(g\) on \(F\) by

\begin{equation*} g(x) = \lim_{k \to \infty} g_k(x) . \end{equation*}

Prove that \(g\) is continuous on \(F\) if and only if \(g_1, g_2, \dots\) converges uniformly on \(F\) to \(g\).

[The result above is called Dini’s Theorem.]

Solution

First note that \(g\) is well defined and real-valued: for each \(x \in F\) the sequence \(g_1(x) \le g_2(x) \le \cdots\) is increasing and bounded above (its supremum is finite by hypothesis), so it converges to

\begin{equation*} g(x) = \sup\{g_1(x), g_2(x), \dots\} \in \mathbb{R} . \end{equation*}

In particular \(g_k(x) \le g(x)\) for every \(k \in \mathbb{Z}^+\) and every \(x \in F\).

If \(F = \emptyset\) both statements hold vacuously, so assume \(F \neq \emptyset\).

One direction requires neither monotonicity nor any hypothesis on \(F\). Suppose \(g_1, g_2, \dots\) converges uniformly on \(F\) to \(g\). Each \(g_k\) is continuous at each point \(b \in F\), so 2.84 (uniform limit of continuous functions is continuous) shows that \(g\) is continuous at each \(b \in F\). Hence \(g\) is continuous on \(F\).

For the other direction, suppose \(g\) is continuous on \(F\). For each \(k \in \mathbb{Z}^+\) define

\begin{equation*} h_k = g - g_k . \end{equation*}

Then each \(h_k\) is a continuous real-valued function on \(F\) (a difference of continuous functions), and

  • \(h_k(x) \ge 0\) for all \(x \in F\), because \(g_k(x) \le g(x)\);
  • \(h_1(x) \ge h_2(x) \ge \cdots\) for all \(x \in F\), because \(g_k(x) \le g_{k+1}(x)\);
  • \(\lim_{k \to \infty} h_k(x) = 0\) for each \(x \in F\), by the definition of \(g\).

Proving that \(g_1, g_2, \dots\) converges uniformly to \(g\) on \(F\) is the same as proving that for every \(\varepsilon > 0\) there exists \(n \in \mathbb{Z}^+\) with \(h_k(x) < \varepsilon\) for all integers \(k \ge n\) and all \(x \in F\); and by the monotonicity of \((h_k)\) in \(k\) it suffices to find, for each \(\varepsilon > 0\), a single \(n\) with \(h_n(x) < \varepsilon\) for all \(x \in F\), since then \(0 \le h_k(x) \le h_n(x) < \varepsilon\) for all \(k \ge n\).

Suppose, seeking a contradiction, that this fails for some \(\varepsilon > 0\). Then for every \(n \in \mathbb{Z}^+\) there exists \(x_n \in F\) with

\begin{equation*} h_n(x_n) \ge \varepsilon . \end{equation*}

The sequence \(x_1, x_2, \dots\) lies in the bounded set \(F\), so by the Bolzano-Weierstrass Theorem it has a convergent subsequence \(x_{n_1}, x_{n_2}, \dots\) with \(n_1 < n_2 < \cdots\), say \(\lim_{j \to \infty} x_{n_j} = x\). Because \(F\) is closed and each \(x_{n_j} \in F\), we have \(x \in F\).

Now fix \(m \in \mathbb{Z}^+\). For every \(j\) large enough that \(n_j \ge m\), monotonicity in the index gives

\begin{equation*} h_m(x_{n_j}) \ge h_{n_j}(x_{n_j}) \ge \varepsilon . \end{equation*}

Since \(h_m\) is continuous at \(x\) and \(x_{n_j} \to x\) with all \(x_{n_j} \in F\), letting \(j \to \infty\) yields

\begin{equation*} h_m(x) = \lim_{j \to \infty} h_m(x_{n_j}) \ge \varepsilon . \end{equation*}

(Continuity of \(h_m\) at \(x\) relative to \(F\) is exactly what licenses passing to the limit here.)

But \(m \in \mathbb{Z}^+\) was arbitrary, so \(h_m(x) \ge \varepsilon\) for every \(m\), contradicting \(\lim_{m \to \infty} h_m(x) = 0\). This contradiction shows that for every \(\varepsilon > 0\) there is an \(n \in \mathbb{Z}^+\) with \(h_n(x) < \varepsilon\) for all \(x \in F\), and hence, as observed above, \(h_k(x) < \varepsilon\) for all \(k \ge n\) and all \(x \in F\). Therefore \(g_1, g_2, \dots\) converges uniformly on \(F\) to \(g\).

Each of the remaining hypotheses is genuinely used in the direction just proved.

Monotonicity of the sequence cannot be dropped. On \(F = [0,1]\) let \(g_k\) be the continuous piecewise-linear function with \(g_k(0) = 0\), \(g_k\bigl(\frac{1}{2k}\bigr) = 1\), \(g_k\bigl(\frac1k\bigr) = 0\), and \(g_k(x) = 0\) for \(x \ge \frac1k\). Then \(g_1, g_2, \dots\) converges pointwise on \([0,1]\) to the continuous function \(g \equiv 0\), yet \(\sup_{x \in [0,1]} |g_k(x) - g(x)| = 1\) for every \(k\), so the convergence is not uniform. This sequence is of course not increasing.

Closedness of \(F\) cannot be dropped. On \(F = (0,1]\), which is bounded but not closed, the functions \(g_k(x) = \min\{kx, 1\}\) are continuous, increasing in \(k\), and bounded above by \(1\); they converge pointwise to the continuous function \(g \equiv 1\). But \(g_k\bigl(\frac{1}{2k}\bigr) = \frac12\), so \(\sup_{x \in (0,1]} |g(x) - g_k(x)| \ge \frac12\) for every \(k\) and the convergence is not uniform. Boundedness cannot be dropped either: on the closed but unbounded set \(F = [0, \infty)\) the functions \(g_k(x) = \max\{1 - \frac{x}{k},\, 0\}\) are continuous, increasing in \(k\), and bounded above by \(1\), and they converge pointwise to \(g \equiv 1\); but \(g_k(x) = 0\) for \(x \ge k\), so \(\sup_{x \ge 0} |g(x) - g_k(x)| = 1\) for every \(k\).

Finally, note that continuity of \(g\) is a genuine hypothesis rather than an automatic consequence of the other assumptions: on \(F = [0,1]\) the functions \(g_k(x) = 1 - x^k\) are continuous, satisfy \(g_1(x) \le g_2(x) \le \cdots\) (since \(x^k\) is nonincreasing in \(k\) for \(x \in [0,1]\)), and have \(\sup\{g_1(x), g_2(x), \dots\} \le 1 < \infty\) for each \(x\); yet their pointwise limit equals \(1\) on \([0,1)\) and \(0\) at \(x = 1\), so \(g\) is not continuous. By the theorem just proved, the convergence is therefore not uniform, as one checks directly from \(g_k\bigl(2^{-1/k}\bigr) = \tfrac12\) while \(g\bigl(2^{-1/k}\bigr) = 1\).

Problem (2E.8)

Suppose \(\mu\) is the measure on \((\mathbb{Z}^+, 2^{\mathbb{Z}^+})\) defined by

\begin{equation*} \mu(E) = \sum_{n \in E} \frac{1}{2^n}. \end{equation*}

Prove that for every \(\varepsilon > 0\), there exists a set \(E \subseteq \mathbb{Z}^+\) with \(\mu(\mathbb{Z}^+ \setminus E) < \varepsilon\) such that \(f_1, f_2, \ldots\) converges uniformly on \(E\) for every sequence of functions \(f_1, f_2, \ldots\) from \(\mathbb{Z}^+\) to \(\mathbb{R}\) that converges pointwise on \(\mathbb{Z}^+\).

[This result does not follow from Egorov’s Theorem because here we are asking for \(E\) to depend only on \(\varepsilon\). In Egorov’s Theorem, \(E\) depends on \(\varepsilon\) and on the sequence \(f_1, f_2, \ldots\).]

Solution

The point is that \(\mu\) concentrates almost all of its mass on a finite set, and pointwise convergence on a finite set is automatically uniform.

Let \(\varepsilon > 0\). Because \(2^n \to \infty\), there exists \(N \in \mathbb{Z}^+\) with \(\frac{1}{2^N} < \varepsilon\). Put

\begin{equation*} E = \{1, 2, \ldots, N\}. \end{equation*}

Then \(\mathbb{Z}^+ \setminus E = \{N+1, N+2, \ldots\}\), and summing the geometric series gives

\begin{equation*} \mu(\mathbb{Z}^+ \setminus E) = \sum_{n=N+1}^{\infty} \frac{1}{2^n} = \frac{1}{2^N} < \varepsilon. \end{equation*}

Note that \(E\) was chosen using only \(\varepsilon\); no sequence of functions was involved.

Now suppose \(f_1, f_2, \ldots\) is any sequence of functions from \(\mathbb{Z}^+\) to \(\mathbb{R}\) that converges pointwise on \(\mathbb{Z}^+\), say to \(f : \mathbb{Z}^+ \to \mathbb{R}\). We show \(f_k \to f\) uniformly on \(E\).

Let \(\delta > 0\). For each \(n \in \{1, \ldots, N\}\), pointwise convergence at \(n\) provides \(M_n \in \mathbb{Z}^+\) such that

\begin{equation*} |f_k(n) - f(n)| < \delta \quad \text{for all } k \geq M_n . \end{equation*}

Let \(M = \max\{M_1, \ldots, M_N\}\), which is a well-defined element of \(\mathbb{Z}^+\) because the maximum is taken over a finite set. If \(k \geq M\) and \(n \in E\), then \(k \geq M_n\) and hence \(|f_k(n) - f(n)| < \delta\). Thus

\begin{equation*} \sup_{n \in E} |f_k(n) - f(n)| \leq \delta \quad \text{for all } k \geq M . \end{equation*}

Because \(\delta > 0\) was arbitrary, \(f_1, f_2, \ldots\) converges uniformly to \(f\) on \(E\).

Hence for each \(\varepsilon > 0\) the single set \(E = \{1, \ldots, N\}\) works simultaneously for every pointwise convergent sequence, as desired.

Problem (2E.9)

Suppose \(F_1, \ldots, F_n\) are disjoint closed subsets of \(\mathbb{R}\). Prove that if

\begin{equation*} g : F_1 \cup \cdots \cup F_n \to \mathbb{R} \end{equation*}

is a function such that \(g|_{F_k}\) is a continuous function for each \(k \in \{1, \ldots, n\}\), then \(g\) is a continuous function.

Solution

Write \(F = F_1 \cup \cdots \cup F_n\). Continuity of \(g\) means continuity of \(g\) as a function on its domain \(F\): for each \(x \in F\) and each \(\varepsilon > 0\) there is \(\delta > 0\) such that \(|g(y) - g(x)| < \varepsilon\) for all \(y \in F\) with \(|y - x| < \delta\).

Fix \(x \in F\) and \(\varepsilon > 0\). Because \(F_1, \ldots, F_n\) cover \(F\), there is \(j \in \{1, \ldots, n\}\) with \(x \in F_j\). Let

\begin{equation*} C = \bigcup_{k \neq j} F_k . \end{equation*}

Being a finite union of closed sets, \(C\) is closed. Because the sets \(F_1, \ldots, F_n\) are disjoint and \(x \in F_j\), we have \(x \notin C\). Hence \(\mathbb{R} \setminus C\) is an open set containing \(x\), so there exists \(\delta_1 > 0\) such that

\begin{equation*} (x - \delta_1, x + \delta_1) \cap C = \varnothing . \end{equation*}

(If \(n = 1\) then \(C = \varnothing\) and any \(\delta_1 > 0\) works.)

Because \(g|_{F_j}\) is continuous at \(x\), there exists \(\delta_2 > 0\) such that

\begin{equation*} |g(y) - g(x)| < \varepsilon \quad \text{for all } y \in F_j \text{ with } |y - x| < \delta_2 . \end{equation*}

Let \(\delta = \min\{\delta_1, \delta_2\} > 0\). Suppose \(y \in F\) and \(|y - x| < \delta\). Then \(|y - x| < \delta_1\), so \(y \notin C\); since \(y \in F = F_j \cup C\), this forces \(y \in F_j\). Also \(|y - x| < \delta_2\), so \(|g(y) - g(x)| < \varepsilon\).

Thus \(g\) is continuous at \(x\). As \(x \in F\) was arbitrary, \(g\) is a continuous function.

Remarks on the hypotheses. Disjointness made the argument easy but is not actually needed. Suppose \(F_1, \ldots, F_n\) are closed but possibly overlapping and \(g|_{F_k}\) is continuous for each \(k\) (which forces the restrictions to agree on overlaps, since they are all restrictions of the single function \(g\)). Fix \(x \in F\) and \(\varepsilon > 0\), let \(P = \{k : x \in F_k\}\), which is nonempty, and let \(C = \bigcup_{k \notin P} F_k\), a closed set with \(x \notin C\); choose \(\delta_1 > 0\) with \((x - \delta_1, x + \delta_1) \cap C = \varnothing\). For each \(k \in P\), continuity of \(g|_{F_k}\) at \(x\) gives \(\eta_k > 0\) such that \(|g(y) - g(x)| < \varepsilon\) whenever \(y \in F_k\) and \(|y - x| < \eta_k\). Let \(\delta = \min\bigl(\{\delta_1\} \cup \{\eta_k : k \in P\}\bigr) > 0\), a minimum over a finite set. If \(y \in F\) and \(|y - x| < \delta\), then \(y \notin C\), so \(y \in F_k\) for some \(k \in P\), and \(|y - x| < \eta_k\) gives \(|g(y) - g(x)| < \varepsilon\). So \(g\) is continuous. This overlapping version is what is usually called the pasting lemma for finitely many closed sets. Closedness, however, is essential: with \(F_1 = (-\infty, 0)\) and \(F_2 = [0, \infty)\) the function \(g = \chi_{F_2}\) restricts to a continuous function on each piece yet is not continuous on \(\mathbb{R}\). Finiteness is essential too: let \(q_1, q_2, \ldots\) enumerate \(\mathbb{Q}\) and take \(F_k = \{q_k\}\), a disjoint sequence of closed sets with union \(\mathbb{Q}\). Every restriction \(g|_{F_k}\) is continuous, because a function on a one-point set is trivially continuous. But the function \(g : \mathbb{Q} \to \mathbb{R}\) defined by \(g(q_1) = 1\) and \(g(q) = 0\) for \(q \neq q_1\) is not continuous at \(q_1\): every interval about \(q_1\) contains rationals other than \(q_1\), at which \(g\) takes the value \(0\).

Problem (2E.10)

Suppose \(F \subseteq \mathbb{R}\) is such that every continuous function from \(F\) to \(\mathbb{R}\) can be extended to a continuous function from \(\mathbb{R}\) to \(\mathbb{R}\). Prove that \(F\) is a closed subset of \(\mathbb{R}\).

Solution

We prove the contrapositive: if \(F\) is not closed, then some continuous function from \(F\) to \(\mathbb{R}\) has no continuous extension to \(\mathbb{R}\).

So suppose \(F\) is not a closed subset of \(\mathbb{R}\). Then \(F\) does not contain all of its limit points, so there exists \(b \in \mathbb{R}\) such that \(b \notin F\) but \(b\) is a limit point of \(F\); equivalently, there is a sequence \(x_1, x_2, \ldots\) of elements of \(F\) with \(\lim_{j \to \infty} x_j = b\) and \(x_j \neq b\) for every \(j\).

Define \(g : F \to \mathbb{R}\) by

\begin{equation*} g(x) = \frac{1}{x - b}. \end{equation*}

This makes sense because \(b \notin F\), so \(x - b \neq 0\) for every \(x \in F\). The function \(x \mapsto \frac{1}{x-b}\) is continuous on \(\mathbb{R} \setminus \{b\}\), and \(F \subseteq \mathbb{R} \setminus \{b\}\); hence \(g\), being the restriction of a continuous function to a subset of its domain, is a continuous function from \(F\) to \(\mathbb{R}\).

Suppose, for contradiction, that there is a continuous function \(h : \mathbb{R} \to \mathbb{R}\) with \(h|_F = g\). Continuity of \(h\) at \(b\) and \(x_j \to b\) give

\begin{equation*} \lim_{j \to \infty} h(x_j) = h(b) \in \mathbb{R}. \end{equation*}

In particular the sequence \(h(x_1), h(x_2), \ldots\) is bounded, being convergent in \(\mathbb{R}\). But \(x_j \in F\), so

\begin{equation*} |h(x_j)| = |g(x_j)| = \frac{1}{|x_j - b|} \to \infty \end{equation*}

as \(j \to \infty\), because \(0 < |x_j - b| \to 0\). Thus \(h(x_1), h(x_2), \ldots\) is unbounded, a contradiction.

Hence \(g\) admits no continuous extension to \(\mathbb{R}\), contradicting the hypothesis on \(F\). Therefore \(F\) is closed.

The converse is also true (though not asked): by the Tietze extension theorem, every continuous real-valued function on a closed subset of \(\mathbb{R}\) extends continuously to \(\mathbb{R}\). For \(F \subseteq \mathbb{R}\) closed this is elementary, since \(\mathbb{R} \setminus F\) is a countable disjoint union of open intervals and one simply extends affinely across each bounded gap and constantly across each unbounded one.

Problem (2E.11)

Prove or give a counterexample: If \(F \subseteq \mathbb{R}\) is such that every bounded continuous function from \(F\) to \(\mathbb{R}\) can be extended to a continuous function from \(\mathbb{R}\) to \(\mathbb{R}\), then \(F\) is a closed subset of \(\mathbb{R}\).

Solution

The statement is true. Restricting attention to bounded functions costs nothing: the unbounded function \(x \mapsto \frac{1}{x-b}\) used in Exercise 10 can be replaced by a bounded function that oscillates rather than blows up.

We prove the contrapositive: if \(F\) is not closed, then some bounded continuous function from \(F\) to \(\mathbb{R}\) has no continuous extension to \(\mathbb{R}\).

Suppose \(F\) is not closed. As in Exercise 10, there is a point \(b \notin F\) that is a limit point of \(F\); because \(b \notin F\), this means \(F \cap (b - \delta, b + \delta) \neq \varnothing\) for every \(\delta > 0\).

We first fix a side. We claim that either \(F \cap (b, b + \delta) \neq \varnothing\) for every \(\delta > 0\), or \(F \cap (b - \delta, b) \neq \varnothing\) for every \(\delta > 0\). Indeed, if both alternatives failed, there would exist \(\delta_1 > 0\) with \(F \cap (b, b + \delta_1) = \varnothing\) and \(\delta_2 > 0\) with \(F \cap (b - \delta_2, b) = \varnothing\); then \(\delta = \min\{\delta_1, \delta_2\}\) would satisfy \(F \cap (b - \delta, b + \delta) = \varnothing\) (using \(b \notin F\) to discard the point \(b\) itself), contradicting the previous paragraph.

Suppose the second alternative holds. Let \(-F = \{-x : x \in F\}\). The map \(x \mapsto -x\) is a homeomorphism of \(\mathbb{R}\) onto itself carrying \(F\) onto \(-F\) and \(b\) to \(-b\), so \(-b \notin -F\) and \((-F) \cap (-b, -b + \delta) = -\bigl(F \cap (b - \delta, b)\bigr) \neq \varnothing\) for every \(\delta > 0\): the first alternative holds for \(-F\) at \(-b\). Moreover, if \(g\) is a bounded continuous function on \(-F\) admitting no continuous extension to \(\mathbb{R}\), then \(x \mapsto g(-x)\) is a bounded continuous function on \(F\) admitting no continuous extension to \(\mathbb{R}\), since a continuous extension \(h\) of the latter would make \(x \mapsto h(-x)\) a continuous extension of the former. So it suffices to treat the first alternative, and we assume from now on that

\begin{equation*} F \cap (b, b + \delta) \neq \varnothing \quad \text{for every } \delta > 0 . \end{equation*}

Construct recursively a strictly decreasing sequence in \(F\) converging to \(b\): choose \(x_1 \in F \cap (b, b+1)\), and having chosen \(x_j\), choose

\begin{equation*} x_{j+1} \in F \cap \Bigl(b, \ b + \tfrac{x_j - b}{2}\Bigr), \end{equation*}

which is nonempty by the displayed property. Then \(b < x_{j+1} < x_j\) for all \(j\) and \(0 < x_j - b \le 2^{1-j}\), so \(x_j \to b\).

Now set

\begin{equation*} a_j = \begin{cases} 1 & \text{if } j \text{ is odd},\\ 0 & \text{if } j \text{ is even},\end{cases} \end{equation*}

and define \(\varphi : \mathbb{R} \setminus \{b\} \to [0,1]\) by

\begin{equation*} \varphi(x) = \begin{cases} 0 & \text{if } x < b,\\ a_{j+1} + (a_j - a_{j+1}) \dfrac{x - x_{j+1}}{x_j - x_{j+1}} & \text{if } x \in [x_{j+1}, x_j] \text{ for some } j \in \mathbb{Z}^+,\\ 1 & \text{if } x \geq x_1 . \end{cases} \end{equation*}

In words: \(\varphi\) is \(0\) to the left of \(b\), and on \((b, x_1]\) it is the piecewise linear function whose graph connects the points \((x_{j+1}, a_{j+1})\) and \((x_j, a_j)\) by line segments.

This definition is legitimate. First, the intervals \([x_{j+1}, x_j]\), \(j \in \mathbb{Z}^+\), together with \([x_1, \infty)\) cover \((b, \infty)\): given \(x > b\), the set \(\{i \in \mathbb{Z}^+ : x \le x_i\}\) is finite because \(x_i \to b < x\), so either it is empty (and then \(x > x_1\)) or it has a largest element \(j\), in which case \(x_{j+1} < x \le x_j\). Second, consecutive formulas agree where the intervals overlap: at \(x = x_j\) the \(j\)-th formula and the \((j-1)\)-st formula both give \(\varphi(x_j) = a_j\), and at \(x = x_1\) the first formula gives \(a_1 = 1\). So \(\varphi\) is a well-defined function with values in \([0,1]\), and

\begin{equation*} \varphi(x_j) = a_j \quad \text{for every } j \in \mathbb{Z}^+ . \end{equation*}

Also \(\varphi\) is continuous on \(\mathbb{R} \setminus \{b\}\). Indeed, \(\varphi\) is constant on the open set \((-\infty, b)\), hence continuous there. And if \(x > b\), choose \(j\) with \(x_j < x\); on the closed half-line \([x_j, \infty)\) the function \(\varphi\) is built from the finitely many closed pieces \([x_j, x_{j-1}], \ldots, [x_2, x_1], [x_1, \infty)\), on each of which it is affine, with matching values at the shared endpoints, so \(\varphi|_{[x_j, \infty)}\) is continuous (this is the pasting principle for finitely many closed sets discussed in Exercise 9). Since \((x_j, \infty)\) is an open neighborhood of \(x\) on which \(\varphi\) agrees with this continuous function, \(\varphi\) is continuous at \(x\).

Now define \(g = \varphi|_F\). Because \(b \notin F\) we have \(F \subseteq \mathbb{R} \setminus \{b\}\), so \(g\) is a continuous function from \(F\) to \(\mathbb{R}\), and \(0 \le g \le 1\), so \(g\) is bounded.

Suppose \(h : \mathbb{R} \to \mathbb{R}\) were continuous with \(h|_F = g\). Since \(x_j \in F\) and \(x_j \to b\), continuity of \(h\) at \(b\) gives

\begin{equation*} h(b) = \lim_{j \to \infty} h(x_j) = \lim_{j \to \infty} g(x_j) = \lim_{j \to \infty} a_j . \end{equation*}

But \(a_1, a_2, \ldots\) is the sequence \(1, 0, 1, 0, \ldots\), which does not converge. This contradiction shows that \(g\) has no continuous extension to \(\mathbb{R}\).

Therefore, if every bounded continuous function from \(F\) to \(\mathbb{R}\) extends to a continuous function on \(\mathbb{R}\), then \(F\) must be closed.

Problem (2E.12)

Give an example of a Borel measurable function \(f\) from \(\mathbb{R}\) to \(\mathbb{R}\) such that there does not exist a set \(B \subseteq \mathbb{R}\) such that \(|\mathbb{R} \setminus B| = 0\) and \(f|_B\) is a continuous function on \(B\).

Solution

Luzin’s Theorem says a Borel measurable function is continuous when restricted to the complement of a set of small measure, and one might hope to shrink that exceptional set to measure \(0\). This exercise shows that hope is false.

The obvious first guess, \(f = \chi_{\mathbb{Q}}\), does not work: take \(B = \mathbb{R} \setminus \mathbb{Q}\); then \(|\mathbb{R} \setminus B| = |\mathbb{Q}| = 0\) and \(f|_B\) is identically \(0\), hence continuous. What we need is a Borel set \(A\) that, together with its complement, has positive measure in every interval, so that no set of full measure can separate \(A\) from its complement locally.

Step 1: a closed nowhere dense set of positive measure inside any interval.

Claim: if \(J\) is a nonempty bounded open interval, there is a closed set \(K \subseteq J\) with empty interior and \(|K| > 0\).

Write \(J = (c,d)\) and set \(c^{\prime} = c + \frac{d-c}{4}\), \(d^{\prime} = d - \frac{d-c}{4}\), so \([c^{\prime}, d^{\prime}] \subseteq J\) and \(d^{\prime} - c^{\prime} = \frac{d-c}{2} > 0\). Let \(q_1, q_2, \ldots\) be an enumeration of \(\mathbb{Q}\), let \(\varepsilon = \frac{d^{\prime}-c^{\prime}}{2}\), and let

\begin{equation*} U = \bigcup_{j=1}^{\infty} \Bigl( q_j - \frac{\varepsilon}{2^{j+1}}, \ q_j + \frac{\varepsilon}{2^{j+1}} \Bigr). \end{equation*}

Then \(U\) is open, \(\mathbb{Q} \subseteq U\), and by countable subadditivity of outer measure (2.8) together with 2.14,

\begin{equation*} |U| \leq \sum_{j=1}^{\infty} \frac{\varepsilon}{2^{j}} = \varepsilon . \end{equation*}

Let \(K = [c^{\prime}, d^{\prime}] \setminus U\). Then \(K\) is closed (it is the intersection of the closed sets \([c^{\prime},d^{\prime}]\) and \(\mathbb{R} \setminus U\)) and \(K \subseteq J\). Since \(K\) contains no rational number and every nonempty open interval contains a rational number, \(K\) has empty interior. Finally \([c^{\prime},d^{\prime}] \subseteq K \cup U\), so subadditivity and 2.14 give

\begin{equation*} d^{\prime} - c^{\prime} \leq |K| + |U| \leq |K| + \varepsilon = |K| + \frac{d^{\prime}-c^{\prime}}{2}, \end{equation*}

hence \(|K| \geq \frac{d^{\prime}-c^{\prime}}{2} > 0\). This proves the claim.

Step 2: a finite union of closed sets with empty interior has empty interior.

It suffices to do two sets, and then induct. Suppose \(E\) and \(G\) are closed with empty interior and \(V \subseteq E \cup G\) with \(V\) open and nonempty. Then \(V \setminus E\) is open; it is nonempty, since \(V \subseteq E\) would contradict \(E\) having empty interior. But \(V \setminus E \subseteq G\), contradicting \(G\) having empty interior. So \(E \cup G\) has empty interior; being a union of two closed sets it is closed.

Step 3: construction of \(A\).

Let \(I_1, I_2, \ldots\) be an enumeration of all open intervals with rational endpoints (a countable collection). We construct recursively closed sets \(K_1, L_1, K_2, L_2, \ldots\), each with empty interior and positive measure, and all pairwise disjoint, with \(K_n \cup L_n \subseteq I_n\).

Suppose \(n \in \mathbb{Z}^+\) and \(K_1, L_1, \ldots, K_{n-1}, L_{n-1}\) have been chosen (for \(n = 1\) nothing has been chosen). Let

\begin{equation*} C = K_1 \cup L_1 \cup \cdots \cup K_{n-1} \cup L_{n-1}, \end{equation*}

a finite union of closed sets with empty interior, hence closed with empty interior by Step 2. Therefore \(I_n \setminus C = I_n \cap (\mathbb{R} \setminus C)\) is open and nonempty (it is nonempty because otherwise the nonempty open set \(I_n\) would be contained in \(C\)). Pick \(y \in I_n \setminus C\). Since \(I_n \setminus C\) is open, there is \(r > 0\) with \((y - r, y + r) \subseteq I_n \setminus C\); set

\begin{equation*} J = \Bigl( y - \tfrac{r}{2}, \ y \Bigr), \qquad J^{\prime} = \Bigl( y, \ y + \tfrac{r}{2} \Bigr), \end{equation*}

which are disjoint nonempty bounded open intervals contained in \(I_n \setminus C\). By Step 1, choose closed sets with empty interior \(K_n \subseteq J\) and \(L_n \subseteq J^{\prime}\) with \(|K_n| > 0\) and \(|L_n| > 0\). By construction \(K_n\) and \(L_n\) are disjoint from each other and from all previously chosen sets, and they are contained in \(I_n\). The recursion continues, producing a pairwise disjoint family.

Now define

\begin{equation*} A = \bigcup_{n=1}^{\infty} K_n , \qquad f = \chi_A . \end{equation*}

Then \(A\) is Borel, being a countable union of closed sets, and \(L_n \cap A = \varnothing\) for every \(n\) by disjointness. The function \(f\) is Borel measurable: for every Borel set \(D \subseteq \mathbb{R}\), the set \(f^{-1}(D)\) is one of \(\varnothing\), \(A\), \(\mathbb{R} \setminus A\), \(\mathbb{R}\), each of which is Borel.

Step 4: \(f\) has the required property.

First note the key density property: for every nonempty open interval \(I\),

\begin{equation*} |A \cap I| > 0 \quad \text{and} \quad |I \setminus A| > 0 . \end{equation*}

Indeed, \(I\) contains some interval \(I_n\) with rational endpoints, and then \(K_n \subseteq A \cap I\) gives \(|A \cap I| \geq |K_n| > 0\), while \(L_n \subseteq I \setminus A\) gives \(|I \setminus A| \geq |L_n| > 0\) (monotonicity of outer measure).

Suppose now that \(B \subseteq \mathbb{R}\) satisfies \(|\mathbb{R} \setminus B| = 0\); we show \(f|_B\) is not continuous.

Since \(|A| \geq |K_1| > 0\) and \(A \subseteq (A \cap B) \cup (\mathbb{R} \setminus B)\), subadditivity gives

\begin{equation*} 0 < |A| \leq |A \cap B| + |\mathbb{R} \setminus B| = |A \cap B| , \end{equation*}

so \(A \cap B \neq \varnothing\). Choose \(x \in A \cap B\), so \(f(x) = 1\).

Suppose, for contradiction, that \(f|_B\) is continuous at \(x\). Taking \(\varepsilon = \frac{1}{2}\), there exists \(\delta > 0\) such that \(|f(y) - f(x)| < \frac{1}{2}\) for all \(y \in B\) with \(|y - x| < \delta\). Since \(f\) takes only the values \(0\) and \(1\), this forces \(f(y) = 1\), i.e. \(y \in A\), for all such \(y\). In other words, with \(I = (x - \delta, x + \delta)\),

\begin{equation*} B \cap I \subseteq A . \end{equation*}

Consequently

\begin{equation*} I \setminus A \subseteq I \setminus (B \cap I) \subseteq \mathbb{R} \setminus B , \end{equation*}

so \(|I \setminus A| \leq |\mathbb{R} \setminus B| = 0\) by monotonicity. This contradicts \(|I \setminus A| > 0\) from the density property.

Hence \(f|_B\) is not continuous at \(x\), and in particular \(f|_B\) is not a continuous function on \(B\). Since \(B\) was an arbitrary subset of \(\mathbb{R}\) with \(|\mathbb{R} \setminus B| = 0\), the Borel measurable function \(f = \chi_A\) is the required example.

Problem (2E.13)

Prove or give a counterexample: If \(f_t : \mathbb{R} \to \mathbb{R}\) is a Borel measurable function for each \(t \in \mathbb{R}\) and \(f : \mathbb{R} \to (-\infty, \infty]\) is defined by

\begin{equation*} f(x) = \sup\{ f_t(x) : t \in \mathbb{R} \}, \end{equation*}

then \(f\) is a Borel measurable function.

Solution

The statement is false. The corresponding result for a countable family is true (2.53: the pointwise supremum of a sequence of \(\mathcal{S}\)-measurable functions is \(\mathcal{S}\)-measurable), but its proof uses the identity

\begin{equation*} h^{-1}\bigl((a, \infty]\bigr) = \bigcup_{k=1}^{\infty} f_k^{-1}\bigl((a, \infty]\bigr) \quad\text{for } h = \sup_k f_k , \end{equation*}

and a \(\sigma\)-algebra is closed only under countable unions. With an uncountable index set the corresponding union is uncountable, and that is exactly what goes wrong.

Counterexample. By 2.67, there exists a set \(A \subseteq \mathbb{R}\) that is not a Borel set. For each \(t \in \mathbb{R}\) define \(f_t : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} f_t = \begin{cases} \chi_{\{t\}} & \text{if } t \in A,\\[2pt] 0 & \text{if } t \notin A .\end{cases} \end{equation*}

Each \(f_t\) is Borel measurable. Indeed, if \(t \notin A\) then \(f_t\) is the constant function \(0\), and constant functions are Borel measurable. If \(t \in A\) then \(f_t = \chi_{\{t\}}\), and for every Borel set \(D \subseteq \mathbb{R}\) the inverse image \(f_t^{-1}(D)\) is one of \(\varnothing\), \(\{t\}\), \(\mathbb{R} \setminus \{t\}\), \(\mathbb{R}\), each of which is a Borel set (the singleton \(\{t\}\) is closed, hence Borel).

Now compute \(f(x) = \sup\{f_t(x) : t \in \mathbb{R}\}\) for \(x \in \mathbb{R}\).

If \(x \in A\), then taking \(t = x\) gives \(f_x(x) = \chi_{\{x\}}(x) = 1\), and since every \(f_t\) takes values in \(\{0,1\}\) we get \(f(x) = 1\).

If \(x \notin A\), fix any \(t \in \mathbb{R}\). If \(t \notin A\) then \(f_t(x) = 0\). If \(t \in A\), then \(t \neq x\) (because \(x \notin A\)), so \(f_t(x) = \chi_{\{t\}}(x) = 0\). Hence \(f_t(x) = 0\) for every \(t\), and \(f(x) = 0\).

Therefore \(f = \chi_A\). Now \((\frac{1}{2}, \infty]\) is a Borel subset of \([-\infty, \infty]\), because by 2.50 this only requires \((\frac{1}{2}, \infty] \cap \mathbb{R} = (\frac{1}{2}, \infty)\) to be a Borel subset of \(\mathbb{R}\), which it is, being open. Since \(f\) takes only the values \(0\) and \(1\),

\begin{equation*} f^{-1}\bigl( (\tfrac{1}{2}, \infty] \bigr) = \{x \in \mathbb{R} : f(x) > \tfrac{1}{2}\} = A , \end{equation*}

which is not a Borel set. By the definition 2.51 of measurability for functions with values in \([-\infty, \infty]\), applied with \(\mathcal{S} = \mathcal{B}\) the \(\sigma\)-algebra of Borel subsets of \(\mathbb{R}\), this shows that \(f\) is not Borel measurable.

Note that all the functions \(f_t\) here are bounded by \(1\), so \(f\) is even real-valued and bounded; the failure has nothing to do with the value \(\infty\) being allowed, only with the uncountability of the index set.

Problem (2E.14)

Suppose \(b_1, b_2, \ldots\) is a sequence of real numbers. Define \(f : \mathbb{R} \to [0, \infty]\) by

\begin{equation*} f(x) = \begin{cases} \displaystyle \sum_{k=1}^{\infty} \frac{1}{4^k |x - b_k|} & \text{if } x \notin \{b_1, b_2, \ldots\},\\[8pt] \infty & \text{if } x \in \{b_1, b_2, \ldots\}. \end{cases} \end{equation*}

Prove that \(|\{x \in \mathbb{R} : f(x) < 1\}| = \infty\).

[This exercise is a variation of a problem originally considered by Borel. If \(b_1, b_2, \ldots\) contains all the rational numbers, then it is not even obvious that \(\{x \in \mathbb{R} : f(x) < \infty\} \neq \varnothing\).]

Solution

Here \(|\cdot|\) denotes outer measure, so no measurability of \(f\) is needed: everything follows from monotonicity and countable subadditivity of outer measure (2.8) together with the formula for the outer measure of an interval (2.14).

Write

\begin{equation*} E = \{x \in \mathbb{R} : f(x) < 1\}, \qquad N = \mathbb{R} \setminus E = \{x \in \mathbb{R} : f(x) \geq 1\}, \end{equation*}

where the value \(f(x) = \infty\) is counted as \(\geq 1\). The strategy: the exceptional set \(N\) is covered by intervals whose lengths form a convergent series, so \(N\) has finite outer measure; a set whose complement has finite outer measure must have infinite outer measure.

Step 1: covering \(N\) by short intervals.

For each \(k \in \mathbb{Z}^+\) let

\begin{equation*} A_k = \Bigl[ b_k - \frac{1}{2^{k-1}}, \ b_k + \frac{1}{2^{k-1}} \Bigr] = \{x \in \mathbb{R} : |x - b_k| \leq 2^{1-k}\}. \end{equation*}

By 2.14,

\begin{equation*} |A_k| = \frac{2}{2^{k-1}} = \frac{4}{2^{k}} . \end{equation*}

We claim

\begin{equation*} N \subseteq \bigcup_{k=1}^{\infty} A_k . \end{equation*}

Equivalently, we show that \(x \notin \bigcup_{k=1}^{\infty} A_k\) implies \(f(x) < 1\). So suppose \(x \notin A_k\) for every \(k\), which means

\begin{equation*} |x - b_k| > \frac{1}{2^{k-1}} \quad \text{for every } k \in \mathbb{Z}^+ . \end{equation*}

In particular \(|x - b_k| > 0\) for every \(k\), so \(x \notin \{b_1, b_2, \ldots\}\) and hence \(f(x)\) is given by the series. For each \(k\), using \(4^k = 2^{2k}\),

\begin{equation*} \frac{1}{4^k |x - b_k|} < \frac{1}{4^k \cdot 2^{1-k}} = \frac{1}{2^{2k} \cdot 2^{1-k}} = \frac{1}{2^{k+1}} . \end{equation*}

Summing the geometric series,

\begin{equation*} f(x) = \sum_{k=1}^{\infty} \frac{1}{4^k |x - b_k|} \leq \sum_{k=1}^{\infty} \frac{1}{2^{k+1}} = \frac{1}{2} < 1 . \end{equation*}

Thus \(x \in E\), which proves the claimed inclusion. (Note that this argument also shows \(E \neq \varnothing\), and in fact that \(f\) is finite off a set of finite measure, no matter how the \(b_k\) are chosen; that answers the concern raised in the bracketed remark.)

Step 2: \(N\) has finite outer measure.

By monotonicity and countable subadditivity of outer measure (2.8),

\begin{equation*} |N| \leq \Bigl| \bigcup_{k=1}^{\infty} A_k \Bigr| \leq \sum_{k=1}^{\infty} |A_k| = \sum_{k=1}^{\infty} \frac{4}{2^k} = 4 . \end{equation*}

Step 3: conclusion.

Let \(n \in \mathbb{Z}^+\). Since \(E \cup N = \mathbb{R}\), we have \([-n, n] \subseteq (E \cap [-n,n]) \cup N\), so subadditivity and 2.14 give

\begin{equation*} 2n = |[-n,n]| \leq |E \cap [-n,n]| + |N| \leq |E| + 4 , \end{equation*}

using monotonicity in the last step. Hence

\begin{equation*} |E| \geq 2n - 4 \quad \text{for every } n \in \mathbb{Z}^+ , \end{equation*}

and letting \(n \to \infty\) forces \(|E| = \infty\). That is,

\begin{equation*} |\{x \in \mathbb{R} : f(x) < 1\}| = \infty , \end{equation*}

as desired.

Two comments. First, the sequence \(b_1, b_2, \ldots\) was completely arbitrary: it may be dense in \(\mathbb{R}\), may repeat values, and may be unbounded, yet the set where \(f < 1\) always has infinite measure while the set where \(f \geq 1\) always has outer measure at most \(4\). The rapidly decaying weights \(4^{-k}\) do all the work, since the set where the \(k\)-th term alone is at least \(2^{-(k+1)}\) is an interval of length \(4 \cdot 2^{-k}\). Second, the same argument with \(A_k\) replaced by \(\{x : |x - b_k| \le \varepsilon 2^{1-k}\}\) shows \(|\{x : f(x) \geq \varepsilon^{-1}\}| \le 4\varepsilon\) for every \(\varepsilon>0\), hence \(|\{x \in \mathbb{R} : f(x) = \infty\}| = 0\): the series converges for almost every \(x \in \mathbb{R}\).

Problem (2E.15)

Suppose \(B\) is a Borel set and \(f : B \to \mathbb{R}\) is a Lebesgue measurable function. Show that there exists a Borel measurable function \(g : B \to \mathbb{R}\) such that

\begin{equation*} |\{x \in B : g(x) \neq f(x)\}| = 0. \end{equation*}

Solution

The strategy is to extend \(f\) to all of \(\mathbb{R}\), apply 2.95 (every Lebesgue measurable function is almost Borel measurable) to the extension, and then restrict back to \(B\). The hypothesis that \(B\) is a Borel set is exactly what makes the last step work.

Step 1: extend \(f\) to \(\mathbb{R}\).

Define \(\tilde{f} : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} \tilde{f}(x) = \begin{cases} f(x) & \text{if } x \in B, \\ 0 & \text{if } x \in \mathbb{R} \setminus B. \end{cases} \end{equation*}

We claim \(\tilde{f}\) is a Lebesgue measurable function in the sense of 2.94, i.e. \(\tilde{f}^{-1}(A)\) is a Lebesgue measurable subset of \(\mathbb{R}\) for every Borel set \(A \subseteq \mathbb{R}\).

So let \(A \subseteq \mathbb{R}\) be a Borel set. Splitting \(\mathbb{R}\) into \(B\) and \(\mathbb{R} \setminus B\), on which \(\tilde{f}\) is respectively \(f\) and the constant \(0\), gives

\begin{equation*} \tilde{f}^{-1}(A) = \begin{cases} f^{-1}(A) & \text{if } 0 \notin A, \\ f^{-1}(A) \cup (\mathbb{R} \setminus B) & \text{if } 0 \in A. \end{cases} \end{equation*}

Here \(f^{-1}(A)\) means \(\{x \in B : f(x) \in A\}\), which is a Lebesgue measurable subset of \(\mathbb{R}\) because \(f\) is a Lebesgue measurable function and \(A\) is a Borel set. Also \(\mathbb{R} \setminus B\) is a Borel set (the Borel sets form a \(\sigma\)-algebra and \(B\) is Borel), hence is Lebesgue measurable, because every Borel set is Lebesgue measurable (the \(\sigma\)-algebra of Lebesgue measurable sets contains all open sets, hence contains the smallest \(\sigma\)-algebra containing the open sets). In either case \(\tilde{f}^{-1}(A)\) is a union of at most two Lebesgue measurable sets and so is Lebesgue measurable. Thus \(\tilde{f}\) is a Lebesgue measurable function from \(\mathbb{R}\) to \(\mathbb{R}\).

Step 2: apply 2.95.

By 2.95, there exists a Borel measurable function \(\tilde{g} : \mathbb{R} \to \mathbb{R}\) such that

\begin{equation*} |\{x \in \mathbb{R} : \tilde{g}(x) \neq \tilde{f}(x)\}| = 0. \end{equation*}

Step 3: restrict to \(B\).

Define \(g : B \to \mathbb{R}\) by \(g = \tilde{g}|_B\), that is, \(g(x) = \tilde{g}(x)\) for \(x \in B\).

First, \(g\) is Borel measurable. Indeed, let \(A \subseteq \mathbb{R}\) be a Borel set. Then

\begin{equation*} g^{-1}(A) = \{x \in B : \tilde{g}(x) \in A\} = B \cap \tilde{g}^{-1}(A). \end{equation*}

Now \(\tilde{g}^{-1}(A)\) is a Borel subset of \(\mathbb{R}\) because \(\tilde{g}\) is Borel measurable, and \(B\) is a Borel set by hypothesis; hence \(g^{-1}(A)\) is a Borel set, being the intersection of two Borel sets. Thus \(g : B \to \mathbb{R}\) is a Borel measurable function.

Second, the exceptional set is small. If \(x \in B\) and \(g(x) \neq f(x)\), then \(\tilde{g}(x) = g(x) \neq f(x) = \tilde{f}(x)\) (the last equality because \(x \in B\)). Hence

\begin{equation*} \{x \in B : g(x) \neq f(x)\} \subseteq \{x \in \mathbb{R} : \tilde{g}(x) \neq \tilde{f}(x)\}. \end{equation*}

Outer measure preserves order (2.5), so

\begin{equation*} |\{x \in B : g(x) \neq f(x)\}| \le |\{x \in \mathbb{R} : \tilde{g}(x) \neq \tilde{f}(x)\}| = 0. \end{equation*}

Therefore \(|\{x \in B : g(x) \neq f(x)\}| = 0\), as desired.

Remark on the hypothesis that \(B\) is Borel.

The assumption that \(B\) is a Borel set is not merely a convenience: it is necessary. By the definition in 2.94, a function \(g : B \to \mathbb{R}\) is Borel measurable when \(g^{-1}(A)\) is a Borel set for every Borel set \(A \subseteq \mathbb{R}\); taking \(A = \mathbb{R}\) forces \(B = g^{-1}(\mathbb{R})\) to be a Borel set. So if \(B\) were Lebesgue measurable but not Borel, no Borel measurable function on \(B\) would exist at all, and the assertion of the exercise would fail for such a \(B\), because its hypothesis can be met: the identity map \(f(x) = x\) on \(B\) is Lebesgue measurable, since \(f^{-1}(A) = B \cap A\) is Lebesgue measurable for every Borel set \(A\). This is also where the argument above would break down: in Step 3 the set \(g^{-1}(A) = B \cap \tilde{g}^{-1}(A)\) is Borel precisely because \(B\) is.

Remark: a self-contained version of Step 2.

For completeness, here is the argument of 2.95 adapted directly to \(B\), avoiding the extension. Let \(\mathcal{L}_B\) denote the \(\sigma\)-algebra of Lebesgue measurable subsets of \(B\) and \(\mathcal{B}_B\) the \(\sigma\)-algebra of Borel subsets of \(B\); because \(B\) is a Borel set, a subset of \(B\) lies in \(\mathcal{B}_B\) exactly when it is a Borel subset of \(\mathbb{R}\), and a subset of \(B\) lies in \(\mathcal{L}_B\) exactly when it is a Lebesgue measurable subset of \(\mathbb{R}\). As noted after 2.94, \(f\) is an \(\mathcal{L}_B\)-measurable function on \(B\). So by 2.89 applied to the measurable space \((B, \mathcal{L}_B)\) there is a sequence \(f_1, f_2, \ldots\) of simple \(\mathcal{L}_B\)-measurable functions from \(B\) to \(\mathbb{R}\) converging pointwise on \(B\) to \(f\). Fix \(k \in \mathbb{Z}^+\) and write

\begin{equation*} f_k = c_1 \chi_{A_1} + \cdots + c_n \chi_{A_n} \end{equation*}

where \(c_1, \ldots, c_n\) are the distinct values taken by \(f_k\) and \(A_j = f_k^{-1}(\{c_j\})\), so that \(A_1, \ldots, A_n\) are disjoint sets in \(\mathcal{L}_B\) whose union is \(B\). Each \(A_j\) is therefore a Lebesgue measurable subset of \(\mathbb{R}\), so by the equivalence of (a) and (d) in 2.71 there is a Borel set \(B_j \subseteq A_j\) with \(|A_j \setminus B_j| = 0\). Put \(g_k = c_1 \chi_{B_1} + \cdots + c_n \chi_{B_n}\), viewed as a function on \(B\); each \(B_j\) is a Borel subset of \(B\), and \(g_k\) takes only the values \(c_1, \ldots, c_n, 0\), each on a set that is a finite union of the Borel sets \(B_1, \ldots, B_n, B \setminus (B_1 \cup \cdots \cup B_n)\), so \(g_k : B \to \mathbb{R}\) is Borel measurable. Moreover

\begin{equation*} \{x \in B : g_k(x) \neq f_k(x)\} \subseteq (A_1 \setminus B_1) \cup \cdots \cup (A_n \setminus B_n) \end{equation*}

has outer measure \(0\) by subadditivity (2.8). Let

\begin{equation*} E = \Bigl\{x \in B : \lim_{k \to \infty} g_k(x) \text{ exists in } \mathbb{R}\Bigr\}, \end{equation*}

which lies in \(\mathcal{B}_B\) by Exercise 14(b) in Section 2B applied to the measurable space \((B, \mathcal{B}_B)\) and the \(\mathcal{B}_B\)-measurable functions \(g_1, g_2, \ldots\); since \(B\) is Borel, \(E\) is thus a Borel subset of \(\mathbb{R}\). If \(x \in B\) lies outside \(\bigcup_{k=1}^{\infty} \{x \in B : g_k(x) \neq f_k(x)\}\), then \(g_k(x) = f_k(x)\) for every \(k\), so \(\lim_{k \to \infty} g_k(x) = f(x)\) exists and \(x \in E\). Hence

\begin{equation*} B \setminus E \subseteq \bigcup_{k=1}^{\infty} \{x \in B : g_k(x) \neq f_k(x)\}, \end{equation*}

so \(|B \setminus E| = 0\) by countable subadditivity (2.8). Define \(g(x) = \lim_{k \to \infty} (\chi_E g_k)(x)\) for \(x \in B\); the limit exists for every \(x \in B\) (by definition of \(E\) when \(x \in E\), and because \((\chi_E g_k)(x) = 0\) for all \(k\) when \(x \in B \setminus E\)). Each \(\chi_E g_k\) is \(\mathcal{B}_B\)-measurable, being the product of the \(\mathcal{B}_B\)-measurable functions \(\chi_E\) (measurable because \(E \in \mathcal{B}_B\)) and \(g_k\), by 2.46(a); so \(g\) is \(\mathcal{B}_B\)-measurable by 2.48, and hence \(g : B \to \mathbb{R}\) is Borel measurable. Finally

\begin{equation*} \{x \in B : g(x) \neq f(x)\} \subseteq (B \setminus E) \cup \bigcup_{k=1}^{\infty} \{x \in B : g_k(x) \neq f_k(x)\}, \end{equation*}

because for \(x \in E\) with \(g_k(x) = f_k(x)\) for all \(k\) we get \(g(x) = \lim_{k \to \infty} g_k(x) = f(x)\). The right side has outer measure \(0\), so \(|\{x \in B : g(x) \neq f(x)\}| = 0\).

Integration

Exercises 3A

Problem (3A.1)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(f : X \to [0, \infty]\) is an \(\mathcal{S}\)-measurable function such that \(\int f \, d\mu < \infty\). Explain why

\begin{equation*} \inf_E f = 0 \end{equation*}

for each set \(E \in \mathcal{S}\) with \(\mu(E) = \infty\).

Solution

Suppose \(E \in \mathcal{S}\) satisfies \(\mu(E) = \infty\), and suppose toward a contradiction that

\begin{equation*} c := \inf_E f > 0 . \end{equation*}

Note first that \(E \neq \varnothing\), since \(\mu(\varnothing) = 0 \neq \infty\); so the infimum above is an infimum of a nonempty subset of \([0, \infty]\) and \(c \in (0, \infty]\).

Let \(\mathcal{P}\) be the \(\mathcal{S}\)-partition of \(X\) consisting of the two sets \(E\) and \(X \setminus E\). (These are disjoint sets in \(\mathcal{S}\) whose union is \(X\), so \(\mathcal{P}\) is indeed an \(\mathcal{S}\)-partition in the sense of 3.1; if \(E = X\) then the second set is \(\varnothing\), which is harmless.) By the definition 3.2 of the lower Lebesgue sum,

\begin{equation*} \mathcal{L}(f, \mathcal{P}) = \mu(E) \inf_E f + \mu(X \setminus E) \inf_{X \setminus E} f \geq \mu(E) \inf_E f = \infty \cdot c = \infty, \end{equation*}

where the inequality holds because the discarded term is an element of \([0, \infty]\) and hence is nonnegative, and where the last equality holds because \(c > 0\).

By the definition 3.3 of the integral, \(\int f \, d\mu = \sup \{ \mathcal{L}(f, \mathcal{P}^{\prime}) : \mathcal{P}^{\prime} \text{ is an } \mathcal{S}\text{-partition of } X \} \geq \mathcal{L}(f, \mathcal{P}) = \infty\). Thus \(\int f \, d\mu = \infty\), contradicting the hypothesis that \(\int f \, d\mu < \infty\).

This contradiction shows that the assumption \(\inf_E f > 0\) is untenable. Since \(f\) takes values in \([0, \infty]\) we always have \(\inf_E f \geq 0\), so \(\inf_E f = 0\), which is what we wanted to prove.

In words: a set of infinite measure on which \(f\) has a positive lower bound would already contribute an infinite lower Lebesgue sum, so a function with finite integral must come arbitrarily close to \(0\) on every set of infinite measure.

Problem (3A.2)

Suppose \(X\) is a set, \(\mathcal{S}\) is a \(\sigma\)-algebra on \(X\), and \(c \in X\). Define the Dirac measure \(\delta_c\) on \((X, \mathcal{S})\) by

\begin{equation*} \delta_c(E) = \begin{cases} 1 & \text{if } c \in E, \\ 0 & \text{if } c \notin E. \end{cases} \end{equation*}

Prove that if \(f : X \to [0, \infty]\) is \(\mathcal{S}\)-measurable, then \(\int f \, d\delta_c = f( c)\).

[Careful: \(\{c\}\) may not be in \(\mathcal{S}\).]

Solution

Throughout, \(\mathcal{L}(f, \mathcal{P})\) denotes the lower Lebesgue sum of 3.2 taken with respect to the measure \(\delta_c\).

Step 1: we show \(\int f \, d\delta_c \leq f( c)\).

Let \(\mathcal{P}\) be an \(\mathcal{S}\)-partition \(A_1, \ldots, A_m\) of \(X\). Because \(A_1, \ldots, A_m\) are disjoint sets whose union is \(X\), the point \(c\) lies in exactly one of them; say \(c \in A_{j_0}\). Then \(\delta_c(A_{j_0}) = 1\) and \(\delta_c(A_j) = 0\) for every \(j \neq j_0\). Hence, using the convention that \(0 \cdot \infty = 0\),

\begin{equation*} \mathcal{L}(f, \mathcal{P}) = \sum_{j=1}^m \delta_c(A_j) \inf_{A_j} f = \inf_{A_{j_0}} f \leq f( c), \end{equation*}

the last inequality because \(c \in A_{j_0}\). Taking the supremum over all \(\mathcal{S}\)-partitions \(\mathcal{P}\) of \(X\) and using the definition 3.3 of the integral gives \(\int f \, d\delta_c \leq f( c)\).

Note that this half of the argument used nothing about \(f\) beyond its being a function into \([0, \infty]\).

Step 2: we show \(\int f \, d\delta_c \geq f( c)\).

Here is where measurability of \(f\) enters, and here is the point of the warning in the exercise: we cannot simply use the partition \(\{c\}, X \setminus \{c\}\), because \(\{c\}\) need not belong to \(\mathcal{S}\). Instead we use level sets of \(f\), which do belong to \(\mathcal{S}\).

If \(f( c) = 0\) the desired inequality is trivial, since integrals of nonnegative functions are nonnegative. So suppose \(f( c) > 0\) and let \(t \in [0, f( c))\) be arbitrary. Set

\begin{equation*} A = f^{-1}\bigl((t, \infty]\bigr) = \{ x \in X : f(x) > t \}. \end{equation*}

The set \((t, \infty]\) is a Borel subset of \([-\infty, \infty]\), because its intersection with \(\mathbb{R}\) is the open interval \((t, \infty)\) (see 2.50). Since \(f\) is \(\mathcal{S}\)-measurable, the definition 2.51 of measurability gives \(A \in \mathcal{S}\). Also \(c \in A\), since \(f( c) > t\).

Let \(\mathcal{P}\) be the \(\mathcal{S}\)-partition of \(X\) consisting of \(A\) and \(X \setminus A\). Since \(c \in A\), we have \(\delta_c(A) = 1\) and \(\delta_c(X \setminus A) = 0\). Therefore, again using \(0 \cdot \infty = 0\),

\begin{equation*} \mathcal{L}(f, \mathcal{P}) = \delta_c(A) \inf_A f + \delta_c(X \setminus A) \inf_{X \setminus A} f = \inf_A f \geq t, \end{equation*}

because \(f(x) > t\) for every \(x \in A\). Hence \(\int f \, d\delta_c \geq t\).

Since \(t \in [0, f( c))\) was arbitrary, taking the supremum over such \(t\) gives \(\int f \, d\delta_c \geq f( c)\). (When \(f( c) = \infty\) this says \(\int f \, d\delta_c \geq t\) for every \(t \in [0, \infty)\), so the integral equals \(\infty = f( c)\).)

Combining Steps 1 and 2 gives \(\int f \, d\delta_c = f( c)\).

Problem (3A.3)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(f : X \to [0, \infty]\) is an \(\mathcal{S}\)-measurable function. Prove that

\begin{equation*} \int f \, d\mu > 0 \quad \text{if and only if} \quad \mu\bigl(\{ x \in X : f(x) > 0 \}\bigr) > 0 . \end{equation*}

Solution

Write

\begin{equation*} E = \{ x \in X : f(x) > 0 \} = f^{-1}\bigl((0, \infty]\bigr) \quad\text{and}\quad E_k = \{ x \in X : f(x) > \tfrac{1}{k} \} = f^{-1}\bigl((\tfrac{1}{k}, \infty]\bigr) \end{equation*}

for \(k \in \mathbb{Z}^+\). Each of \((0, \infty]\) and \((\frac{1}{k}, \infty]\) is a Borel subset of \([-\infty, \infty]\), since its intersection with \(\mathbb{R}\) is an open interval (see 2.50). Because \(f\) is \(\mathcal{S}\)-measurable, the definition 2.51 gives \(E \in \mathcal{S}\) and \(E_k \in \mathcal{S}\) for each \(k\).

First suppose \(\mu(E) > 0\); we show \(\int f \, d\mu > 0\).

Because \(\frac{1}{1} \geq \frac{1}{2} \geq \cdots\), we have \(E_1 \subseteq E_2 \subseteq \cdots\). Moreover \(\bigcup_{k=1}^{\infty} E_k = E\): the inclusion \(\subseteq\) is clear since \(\frac{1}{k} > 0\); conversely if \(f(x) > 0\) then \(\frac{1}{k} < f(x)\) for some \(k \in \mathbb{Z}^+\), so \(x \in E_k\). Hence 2.59 (measure of an increasing union) gives

\begin{equation*} \lim_{k \to \infty} \mu(E_k) = \mu(E) > 0 . \end{equation*}

Therefore there exists \(k \in \mathbb{Z}^+\) with \(\mu(E_k) > 0\). Fix such a \(k\).

Now \(\frac{1}{k} \chi_{E_k}(x) \leq f(x)\) for every \(x \in X\): if \(x \in E_k\) then \(f(x) > \frac{1}{k}\), and if \(x \notin E_k\) then the left side is \(0 \leq f(x)\). Consequently, by 3.8 (integration is order preserving) and then 3.7 (integral of a simple function, applied with the single set \(E_k\) and the coefficient \(\frac{1}{k}\)),

\begin{equation*} \int f \, d\mu \geq \int \tfrac{1}{k} \chi_{E_k} \, d\mu = \tfrac{1}{k} \mu(E_k) > 0 , \end{equation*}

as desired.

Conversely suppose \(\mu(E) = 0\); we show \(\int f \, d\mu = 0\), which proves the remaining implication in contrapositive form.

Let \(\mathcal{P}\) be an arbitrary \(\mathcal{S}\)-partition \(A_1, \ldots, A_m\) of \(X\). Fix \(j \in \{1, \ldots, m\}\) and consider the term \(\mu(A_j) \inf_{A_j} f\) of the lower Lebesgue sum. There are two cases.

If \(A_j \subseteq E\), then \(\mu(A_j) \leq \mu(E) = 0\) by 2.57(a), so \(\mu(A_j) = 0\) and the term equals \(0 \cdot \inf_{A_j} f = 0\) (using the convention \(0 \cdot \infty = 0\)).

If \(A_j \not\subseteq E\), then there exists \(x \in A_j\) with \(x \notin E\), which means \(f(x) = 0\) (recall \(f\) takes values in \([0, \infty]\)). Hence \(\inf_{A_j} f = 0\) and the term equals \(\mu(A_j) \cdot 0 = 0\) (using the convention \(\infty \cdot 0 = 0\)).

In either case the term is \(0\), so \(\mathcal{L}(f, \mathcal{P}) = 0\). Since \(\mathcal{P}\) was an arbitrary \(\mathcal{S}\)-partition of \(X\), the definition 3.3 gives

\begin{equation*} \int f \, d\mu = \sup \{ \mathcal{L}(f, \mathcal{P}) : \mathcal{P} \text{ is an } \mathcal{S}\text{-partition of } X \} = 0 . \end{equation*}

Thus \(\int f \, d\mu > 0\) if and only if \(\mu(E) > 0\).

Problem (3A.4)

Give an example of a Borel measurable function \(f : [0, 1] \to (0, \infty)\) such that \(L(f, [0, 1]) = 0\).

[Recall that \(L(f, [0, 1])\) denotes the lower Riemann integral, which was defined in Section 1A. If \(\lambda\) is Lebesgue measure on \([0, 1]\), then the previous exercise states that \(\int f \, d\lambda > 0\) for this function \(f\), which is what we expect of a positive function. Thus even though both \(L(f, [0, 1])\) and \(\int f \, d\lambda\) are defined by taking the supremum of approximations from below, Lebesgue measure captures the right behavior for this function \(f\) and the lower Riemann integral does not.]

Solution

Let \(r_1, r_2, \ldots\) be an enumeration (a bijective listing) of the rational numbers in \([0, 1]\); such a listing exists because \(\mathbb{Q} \cap [0, 1]\) is countably infinite. Define \(f : [0, 1] \to (0, \infty)\) by

\begin{equation*} f(x) = \begin{cases} \dfrac{1}{k} & \text{if } x = r_k \text{ for some } k \in \mathbb{Z}^+, \\[2mm] 1 & \text{if } x \text{ is irrational.} \end{cases} \end{equation*}

Because the listing is bijective, each rational \(x \in [0, 1]\) equals \(r_k\) for exactly one \(k\), so \(f\) is well defined. Every value of \(f\) lies in \((0, 1] \subseteq (0, \infty)\), so \(f\) indeed maps into \((0, \infty)\); note also that \(f\) is bounded, so its lower Riemann integral on \([0, 1]\) is defined.

The function \(f\) is Borel measurable. Indeed, let \(B \subseteq \mathbb{R}\) be a Borel set. Then

\begin{equation*} f^{-1}(B) = \{ r_k : k \in \mathbb{Z}^+ \text{ and } \tfrac{1}{k} \in B \} \; \cup \; \begin{cases} [0,1] \setminus \mathbb{Q} & \text{if } 1 \in B, \\ \varnothing & \text{if } 1 \notin B. \end{cases} \end{equation*}

The first set is countable, hence Borel (a countable set is a countable union of singletons, and singletons are closed, hence Borel). The second set is either \(\varnothing\) or \([0,1] \setminus \mathbb{Q}\), and \([0,1] \setminus \mathbb{Q}\) is the intersection of the Borel set \([0,1]\) with the complement of the countable, hence Borel, set \(\mathbb{Q}\). So \(f^{-1}(B)\) is a union of two Borel sets and is therefore Borel. Thus \(f\) is Borel measurable.

Now we compute the lower Riemann integral. Let \(P\) be an arbitrary partition \(0 = x_0 < x_1 < \cdots < x_n = 1\) of \([0, 1]\), and fix \(j \in \{1, \ldots, n\}\). The interval \([x_{j-1}, x_j]\) is a nondegenerate closed subinterval of \([0, 1]\), so it contains infinitely many rational numbers. Since the listing \(r_1, r_2, \ldots\) is injective, the index set

\begin{equation*} K_j = \{ k \in \mathbb{Z}^+ : r_k \in [x_{j-1}, x_j] \} \end{equation*}

is infinite, and therefore contains arbitrarily large integers. Hence

\begin{equation*} \inf_{[x_{j-1}, x_j]} f \; \leq \; \inf \Bigl\{ \tfrac{1}{k} : k \in K_j \Bigr\} = 0 , \end{equation*}

and since \(f > 0\) this infimum equals \(0\). Therefore, by the definition 1.3 of the lower Riemann sum,

\begin{equation*} L(f, P, [0,1]) = \sum_{j=1}^n (x_j - x_{j-1}) \inf_{[x_{j-1}, x_j]} f = \sum_{j=1}^n (x_j - x_{j-1}) \cdot 0 = 0 . \end{equation*}

Because \(P\) was an arbitrary partition of \([0, 1]\), taking the supremum over all partitions and using the definition 1.7 of the lower Riemann integral gives

\begin{equation*} L(f, [0,1]) = 0 , \end{equation*}

as required.

For contrast, let \(\lambda\) denote Lebesgue measure on the Borel subsets of \([0, 1]\). Since \(f(x) > 0\) for every \(x \in [0,1]\), we have \(\lambda(\{x \in [0,1] : f(x) > 0\}) = \lambda([0,1]) = 1 > 0\), so the previous exercise (3A.3) gives \(\int f \, d\lambda > 0\). In fact the value can be computed exactly: \(\chi_{[0,1] \setminus \mathbb{Q}} \leq f \leq \chi_{[0,1]}\) pointwise on \([0,1]\), so 3.8 and 3.4 give

\begin{equation*} 1 = \lambda\bigl([0,1] \setminus \mathbb{Q}\bigr) \leq \int f \, d\lambda \leq \lambda([0,1]) = 1 , \end{equation*}

so \(\int f \, d\lambda = 1\). The lower Riemann integral fails here because Riemann partitions are forced to use intervals, and every nondegenerate interval contains points where \(f\) is nearly \(0\); Lebesgue’s definition may use arbitrary Borel sets in a partition, and the set \([0,1] \setminus \mathbb{Q}\) (on which \(f\) equals \(1\)) is allowed as a partition piece.

Problem (3A.5)

Verify the assertion that integration with respect to counting measure is summation (Example 3.6).

Solution

Example 3.6 asserts the following. Suppose \(\mu\) is counting measure on \(\mathbb{Z}^+\) and \(b_1, b_2, \ldots\) is a sequence of nonnegative numbers; think of \(b\) as the function from \(\mathbb{Z}^+\) to \([0, \infty)\) defined by \(b(k) = b_k\). Then

\begin{equation*} \int b \, d\mu = \sum_{k=1}^{\infty} b_k . \end{equation*}

Here \(\mathcal{S}\) is the \(\sigma\)-algebra of all subsets of \(\mathbb{Z}^+\) and \(\mu\) is counting measure (2.55): \(\mu(E)\) is the number of elements of \(E\) when \(E\) is finite, and \(\mu(E) = \infty\) when \(E\) is infinite. Every function on \(\mathbb{Z}^+\) is \(\mathcal{S}\)-measurable, since every subset of \(\mathbb{Z}^+\) is in \(\mathcal{S}\). Note that \(\sum_{k=1}^{\infty} b_k\) is a series of nonnegative terms and so converges in \([0, \infty]\), its value being \(\lim_{n \to \infty} \sum_{k=1}^n b_k\); equivalently it is \(\sup \{ \sum_{k \in F} b_k : F \subseteq \mathbb{Z}^+ \text{ finite} \}\).

Both inequalities are proved directly from the definitions 3.2 and 3.3.

First we show \(\int b \, d\mu \geq \sum_{k=1}^{\infty} b_k\).

Fix \(n \in \mathbb{Z}^+\) and let \(\mathcal{P}_n\) be the \(\mathcal{S}\)-partition of \(\mathbb{Z}^+\) consisting of the \(n+1\) sets

\begin{equation*} \{1\}, \{2\}, \ldots, \{n\}, \quad \{n+1, n+2, \ldots\} . \end{equation*}

These are disjoint and their union is \(\mathbb{Z}^+\). Since \(\mu(\{k\}) = 1\) and \(\inf_{\{k\}} b = b_k\), and since the last term of the lower Lebesgue sum is an element of \([0, \infty]\) and hence nonnegative,

\begin{equation*} \mathcal{L}(b, \mathcal{P}_n) = \sum_{k=1}^n b_k + \mu(\{n+1, n+2, \ldots\}) \inf_{k > n} b_k \geq \sum_{k=1}^n b_k . \end{equation*}

Hence \(\int b \, d\mu \geq \sum_{k=1}^n b_k\) for every \(n \in \mathbb{Z}^+\). Letting \(n \to \infty\) gives \(\int b \, d\mu \geq \sum_{k=1}^{\infty} b_k\).

Next we show \(\int b \, d\mu \leq \sum_{k=1}^{\infty} b_k\).

Let \(\mathcal{P}\) be an arbitrary \(\mathcal{S}\)-partition \(A_1, \ldots, A_m\) of \(\mathbb{Z}^+\). We claim that for each \(j\),

\begin{equation*} \mu(A_j) \inf_{A_j} b \leq \sum_{k \in A_j} b_k , \end{equation*}

where the right side denotes \(\sup \{ \sum_{k \in F} b_k : F \subseteq A_j \text{ finite} \} \in [0, \infty]\). To see this, put \(c_j = \inf_{A_j} b\) and consider three cases.

If \(A_j = \varnothing\), then the left side is \(0 \cdot \infty = 0\) and the right side is \(0\).

If \(A_j\) is finite and nonempty, say with \(n_j\) elements, then \(\mu(A_j) = n_j\) and each of the \(n_j\) numbers \(b_k\) with \(k \in A_j\) is at least \(c_j\), so \(\mu(A_j) c_j = n_j c_j \leq \sum_{k \in A_j} b_k\).

If \(A_j\) is infinite, then \(\mu(A_j) = \infty\). If \(c_j = 0\) the left side is \(\infty \cdot 0 = 0\) and the inequality is clear. If \(c_j > 0\), then for every \(N \in \mathbb{Z}^+\) we may choose a subset \(F \subseteq A_j\) with \(N\) elements, giving \(\sum_{k \in A_j} b_k \geq \sum_{k \in F} b_k \geq N c_j\); letting \(N \to \infty\) shows \(\sum_{k \in A_j} b_k = \infty\), so the inequality holds.

This proves the claim. Summing the claim over \(j = 1, \ldots, m\) gives

\begin{equation*} \mathcal{L}(b, \mathcal{P}) = \sum_{j=1}^m \mu(A_j) \inf_{A_j} b \leq \sum_{j=1}^m \sum_{k \in A_j} b_k \leq \sum_{k=1}^{\infty} b_k . \end{equation*}

To justify the last inequality, write \(s_j = \sum_{k \in A_j} b_k \in [0, \infty]\) and let \(t_1, \ldots, t_m \in [0, \infty)\) be any real numbers with \(t_j < s_j\) whenever \(s_j > 0\) and \(t_j = 0\) whenever \(s_j = 0\). By the definition of \(s_j\) as a supremum, for each \(j\) there is a finite set \(F_j \subseteq A_j\) with \(\sum_{k \in F_j} b_k \geq t_j\). The sets \(F_1, \ldots, F_m\) are disjoint (being subsets of the disjoint sets \(A_1, \ldots, A_m\)) and their union is a finite subset of \(\mathbb{Z}^+\), so

\begin{equation*} \sum_{j=1}^m t_j \; \leq \; \sum_{j=1}^m \sum_{k \in F_j} b_k \; = \sum_{k \in F_1 \cup \cdots \cup F_m} b_k \; \leq \; \sum_{k=1}^{\infty} b_k . \end{equation*}

Now take the supremum over all such tuples \((t_1, \ldots, t_m)\). That supremum equals \(\sum_{j=1}^m s_j\): if every \(s_j\) is finite the \(t_j\) may be taken arbitrarily close to \(s_j\), while if some \(s_j = \infty\) the corresponding \(t_j\) may be taken arbitrarily large, making both sides \(\infty\). Hence \(\sum_{j=1}^m \sum_{k \in A_j} b_k \leq \sum_{k=1}^{\infty} b_k\).

Since \(\mathcal{P}\) was arbitrary, the definition 3.3 gives \(\int b \, d\mu \leq \sum_{k=1}^{\infty} b_k\).

Combining the two inequalities gives \(\int b \, d\mu = \sum_{k=1}^{\infty} b_k\), verifying Example 3.6.

Alternatively, once the Monotone Convergence Theorem 3.11 is available there is a two-line proof. For \(n \in \mathbb{Z}^+\) define \(b^{(n)} = \sum_{k=1}^n b_k \chi_{\{k\}}\). Then \(0 \leq b^{(1)} \leq b^{(2)} \leq \cdots\) (because each \(b_k \geq 0\)) and \(\lim_{n \to \infty} b^{(n)}(k) = b_k\) for every \(k \in \mathbb{Z}^+\) (indeed \(b^{(n)}(k) = b_k\) once \(n \geq k\)). Hence 3.11 and 3.7 give

\begin{equation*} \int b \, d\mu = \lim_{n \to \infty} \int b^{(n)} \, d\mu = \lim_{n \to \infty} \sum_{k=1}^n b_k \mu(\{k\}) = \lim_{n \to \infty} \sum_{k=1}^n b_k = \sum_{k=1}^{\infty} b_k . \end{equation*}

Problem (3A.6)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space, \(f : X \to [0, \infty]\) is \(\mathcal{S}\)-measurable, and \(\mathcal{P}\) and \(\mathcal{P}^{\prime}\) are \(\mathcal{S}\)-partitions of \(X\) such that each set in \(\mathcal{P}^{\prime}\) is contained in some set in \(\mathcal{P}\). Prove that \(\mathcal{L}(f, \mathcal{P}) \leq \mathcal{L}(f, \mathcal{P}^{\prime})\).

Solution

This is the analogue for lower Lebesgue sums of 1.5, which says that refining a partition increases the lower Riemann sum.

Write \(\mathcal{P}\) as the list \(A_1, \ldots, A_m\) and \(\mathcal{P}^{\prime}\) as the list \(B_1, \ldots, B_n\). By hypothesis, for each \(k \in \{1, \ldots, n\}\) there exists an index \(j(k) \in \{1, \ldots, m\}\) with

\begin{equation*} B_k \subseteq A_{j(k)} . \end{equation*}

Fix such a choice of \(j(1), \ldots, j(n)\) once and for all (if some \(B_k\) is contained in more than one \(A_j\), which can happen only when \(B_k = \varnothing\), choose either).

We first record a set-theoretic fact.

Claim: for each \(j \in \{1, \ldots, m\}\),

\begin{equation*} A_j = \bigcup_{\{k \, : \, j(k) = j\}} B_k , \end{equation*}

and the sets in this union are disjoint.

To prove the claim, note first that disjointness is inherited from \(\mathcal{P}^{\prime}\). The inclusion \(\supseteq\) is immediate, since \(j(k) = j\) gives \(B_k \subseteq A_j\). For \(\subseteq\), let \(x \in A_j\). Because \(\mathcal{P}^{\prime}\) is a partition of \(X\), there is a (unique) \(k\) with \(x \in B_k\). Then \(x \in B_k \subseteq A_{j(k)}\), so \(x \in A_j \cap A_{j(k)}\). Since \(A_1, \ldots, A_m\) are disjoint and this intersection is nonempty, \(j(k) = j\). Thus \(x\) lies in the union on the right, proving the claim.

By finite additivity of \(\mu\) (a consequence of the definition 2.54 of a measure), the claim gives

\begin{equation*} \mu(A_j) = \sum_{\{k \, : \, j(k) = j\}} \mu(B_k) \qquad \text{for each } j \in \{1, \ldots, m\}. \end{equation*}

We also note the elementary distributive law in \([0, \infty]\) that will be used: if \(c \in [0, \infty]\) and \(t_1, \ldots, t_p \in [0, \infty]\), then \(\sum_{i=1}^p c \, t_i = c \sum_{i=1}^p t_i\), where \(0 \cdot \infty = \infty \cdot 0 = 0\). Indeed, for \(c \in (0, \infty)\) this is the usual distributive law (both sides being \(\infty\) exactly when some \(t_i = \infty\)); for \(c = 0\) both sides equal \(0\); and for \(c = \infty\) both sides equal \(0\) if every \(t_i = 0\) and equal \(\infty\) otherwise.

Now we compute. The index sets \(\{k : j(k) = j\}\), for \(j = 1, \ldots, m\), are disjoint and their union is \(\{1, \ldots, n\}\), so we may group the terms of \(\mathcal{L}(f, \mathcal{P}^{\prime})\) according to \(j\):

\begin{equation*} \begin{aligned} \mathcal{L}(f, \mathcal{P}^{\prime}) &= \sum_{k=1}^n \mu(B_k) \inf_{B_k} f \\ &= \sum_{j=1}^m \; \sum_{\{k \, : \, j(k) = j\}} \mu(B_k) \inf_{B_k} f \\ &\geq \sum_{j=1}^m \; \sum_{\{k \, : \, j(k) = j\}} \mu(B_k) \inf_{A_j} f \\ &= \sum_{j=1}^m \Bigl( \inf_{A_j} f \Bigr) \sum_{\{k \, : \, j(k) = j\}} \mu(B_k) \\ &= \sum_{j=1}^m \mu(A_j) \inf_{A_j} f \\ &= \mathcal{L}(f, \mathcal{P}) . \end{aligned} \end{equation*}

The inequality in the third line holds because \(B_k \subseteq A_j\) whenever \(j(k) = j\), and the infimum of \(f\) over a subset is at least the infimum over the larger set, so \(\inf_{B_k} f \geq \inf_{A_j} f\); multiplying by \(\mu(B_k) \in [0, \infty]\) preserves this inequality (with the convention \(0 \cdot \infty = 0\), since if \(\mu(B_k) = 0\) both sides are \(0\)). The fourth line uses the distributive law recorded above, and the fifth line uses finite additivity as recorded above.

Hence \(\mathcal{L}(f, \mathcal{P}) \leq \mathcal{L}(f, \mathcal{P}^{\prime})\), as desired.

Note the degenerate cases cause no trouble: if \(B_k = \varnothing\) then \(\inf_{B_k} f = \infty\) but \(\mu(B_k) = 0\), so that term contributes \(0\) on both sides; similarly for empty \(A_j\).

Problem (3A.7)

Suppose \(X\) is a set, \(\mathcal{S}\) is the \(\sigma\)-algebra of all subsets of \(X\), and \(w : X \to [0, \infty]\) is a function. Define a measure \(\mu\) on \((X, \mathcal{S})\) by

\begin{equation*} \mu(E) = \sum_{x \in E} w(x) \end{equation*}

for \(E \subseteq X\). Prove that if \(f : X \to [0, \infty]\) is a function, then

\begin{equation*} \int f \, d\mu = \sum_{x \in X} w(x) f(x), \end{equation*}

where the infinite sums above are defined as the supremum of all sums over finite subsets of \(E\) (first sum) or \(X\) (second sum).

Solution

Notation. For a function \(g : X \to [0, \infty]\) and a set \(A \subseteq X\), write

\begin{equation*} \sum_{x \in A} g(x) = \sup \Bigl\{ \sum_{x \in F} g(x) : F \subseteq A \text{ and } F \text{ is finite} \Bigr\} \in [0, \infty] , \end{equation*}

with the convention that the empty sum is \(0\) (so the supremum is over a nonempty collection and is well defined, and equals \(0\) when \(A = \varnothing\)). Throughout, \(0 \cdot \infty = \infty \cdot 0 = 0\).

Note that every function \(f : X \to [0, \infty]\) is \(\mathcal{S}\)-measurable, because \(\mathcal{S}\) consists of all subsets of \(X\); so the integral \(\int f \, d\mu\) is defined.

Preliminary 0 (a supremum lemma, used repeatedly). Suppose \(A_1, \ldots, A_m\) are disjoint subsets of \(X\) with union contained in a set \(A \subseteq X\), and \(g : X \to [0, \infty]\). Then

\begin{equation*} \sum_{j=1}^m \; \sum_{x \in A_j} g(x) \; \leq \; \sum_{x \in A} g(x) . \end{equation*}

To see this, write \(s_j = \sum_{x \in A_j} g(x)\) and let \(t_1, \ldots, t_m \in [0, \infty)\) satisfy \(t_j < s_j\) when \(s_j > 0\) and \(t_j = 0\) when \(s_j = 0\). By definition of \(s_j\) as a supremum, choose a finite \(F_j \subseteq A_j\) with \(\sum_{x \in F_j} g(x) \geq t_j\). The sets \(F_1, \ldots, F_m\) are disjoint and \(F_1 \cup \cdots \cup F_m\) is a finite subset of \(A\), so

\begin{equation*} \sum_{j=1}^m t_j \leq \sum_{j=1}^m \sum_{x \in F_j} g(x) = \sum_{x \in F_1 \cup \cdots \cup F_m} g(x) \leq \sum_{x \in A} g(x) . \end{equation*}

Taking the supremum over all such tuples \((t_1, \ldots, t_m)\) gives \(\sum_{j=1}^m s_j\) on the left: if every \(s_j\) is finite the \(t_j\) may be taken arbitrarily close to \(s_j\), and if some \(s_j = \infty\) the corresponding \(t_j\) may be taken arbitrarily large, forcing both sides to be \(\infty\). This proves the lemma.

Preliminary 1: \(\mu\) is a measure. (This is asserted without proof in 2.55; here is the verification.) Clearly \(\mu(\varnothing) = 0\), since the only finite subset of \(\varnothing\) is \(\varnothing\) and the empty sum is \(0\). Let \(E_1, E_2, \ldots\) be disjoint subsets of \(X\) with union \(E\). If \(F \subseteq E\) is finite, then \(F\) meets only finitely many \(E_i\), so for some \(n\) we have \(\sum_{x \in F} w(x) = \sum_{i=1}^n \sum_{x \in F \cap E_i} w(x) \leq \sum_{i=1}^n \mu(E_i) \leq \sum_{i=1}^{\infty} \mu(E_i)\), the first inequality because \(F \cap E_i\) is a finite subset of \(E_i\); taking the supremum over such \(F\) gives \(\mu(E) \leq \sum_{i=1}^{\infty} \mu(E_i)\). Conversely, for each \(n\), Preliminary 0 applied to the disjoint sets \(E_1, \ldots, E_n \subseteq E\) and to \(g = w\) gives \(\sum_{i=1}^n \mu(E_i) \leq \mu(E)\); letting \(n \to \infty\) gives \(\sum_{i=1}^{\infty} \mu(E_i) \leq \mu(E)\). Hence \(\mu\) is countably additive and is a measure.

Preliminary 2: \(\mu(\{x\}) = w(x)\) for each \(x \in X\), since the only finite subsets of \(\{x\}\) are \(\varnothing\) and \(\{x\}\), with sums \(0\) and \(w(x)\).

Preliminary 3: suppose \(c \in [0, \infty]\), \(A \subseteq X\), and \(g : X \to [0, \infty]\) satisfies \(g(x) \geq c \, w(x)\) for all \(x \in A\). Then for every finite \(F \subseteq A\) we have \(\sum_{x \in F} g(x) \geq \sum_{x \in F} c \, w(x) = c \sum_{x \in F} w(x)\), the last equality being the distributive law for a finite sum in \([0, \infty]\) (for \(c \in (0, \infty)\) it is standard; for \(c = 0\) both sides are \(0\); for \(c = \infty\) both sides are \(0\) if every \(w(x)\) with \(x \in F\) is \(0\) and are \(\infty\) otherwise). Hence

\begin{equation*} \sum_{x \in A} g(x) \geq \sup_{F} \; c \sum_{x \in F} w(x) = c \sum_{x \in A} w(x) . \end{equation*}

The last equality is the identity \(\sup_F (c \, t_F) = c \sup_F t_F\) for \(c \in [0, \infty]\) and \(t_F \in [0, \infty]\): for \(c \in (0, \infty)\) this is standard; for \(c = 0\) both sides are \(0\); and for \(c = \infty\) both sides are \(0\) if every \(t_F = 0\) and both are \(\infty\) otherwise.

Now we prove the two inequalities. Write \(S = \sum_{x \in X} w(x) f(x)\).

Step 1: \(\int f \, d\mu \geq S\).

Let \(F = \{x_1, \ldots, x_p\}\) be an arbitrary finite subset of \(X\), with \(x_1, \ldots, x_p\) distinct. Let \(\mathcal{P}\) be the \(\mathcal{S}\)-partition of \(X\) consisting of the \(p + 1\) sets

\begin{equation*} \{x_1\}, \ldots, \{x_p\}, \quad X \setminus F , \end{equation*}

which are disjoint subsets of \(X\) (hence in \(\mathcal{S}\)) whose union is \(X\). Using Preliminary 2 and \(\inf_{\{x_i\}} f = f(x_i)\), and discarding the final nonnegative term,

\begin{equation*} \mathcal{L}(f, \mathcal{P}) = \sum_{i=1}^p w(x_i) f(x_i) + \mu(X \setminus F) \inf_{X \setminus F} f \; \geq \; \sum_{x \in F} w(x) f(x) . \end{equation*}

Therefore \(\int f \, d\mu \geq \sum_{x \in F} w(x) f(x)\) by 3.3. Taking the supremum over all finite \(F \subseteq X\) gives \(\int f \, d\mu \geq S\).

Step 2: \(\int f \, d\mu \leq S\).

Let \(\mathcal{P}\) be an arbitrary \(\mathcal{S}\)-partition \(A_1, \ldots, A_m\) of \(X\), and fix \(j\). Put \(c_j = \inf_{A_j} f\). Then \(f(x) \geq c_j\) for every \(x \in A_j\), so \(w(x) f(x) \geq c_j w(x)\) for every \(x \in A_j\) (multiplication by \(w(x) \in [0, \infty]\) preserves the inequality under our conventions: if \(w(x) = 0\) both sides are \(0\)). Applying Preliminary 3 with \(g(x) = w(x) f(x)\), \(A = A_j\), \(c = c_j\) gives

\begin{equation*} \sum_{x \in A_j} w(x) f(x) \; \geq \; c_j \sum_{x \in A_j} w(x) \; = \; c_j \, \mu(A_j) = \mu(A_j) \inf_{A_j} f . \end{equation*}

(If \(A_j = \varnothing\), then \(\mu(A_j) = 0\) and both sides are \(0\).) Summing over \(j = 1, \ldots, m\),

\begin{equation*} \mathcal{L}(f, \mathcal{P}) = \sum_{j=1}^m \mu(A_j) \inf_{A_j} f \; \leq \; \sum_{j=1}^m \; \sum_{x \in A_j} w(x) f(x) \; \leq \; S . \end{equation*}

The last inequality is Preliminary 0 applied to the disjoint sets \(A_1, \ldots, A_m\) (whose union is \(X\)), with \(A = X\) and \(g(x) = w(x) f(x)\).

Since \(\mathcal{P}\) was arbitrary, 3.3 gives \(\int f \, d\mu \leq S\).

Combining Steps 1 and 2,

\begin{equation*} \int f \, d\mu = \sum_{x \in X} w(x) f(x) . \end{equation*}

Taking \(w(x) = 1\) for all \(x\) makes \(\mu\) counting measure on \(X\), and the result then contains Exercise 3A.5 (with \(X = \mathbb{Z}^+\)) as a special case; taking \(w = \chi_{\{c\}}\) makes \(\mu\) the Dirac measure \(\delta_c\) on the \(\sigma\)-algebra of all subsets of \(X\), recovering Exercise 3A.2 in that case.

Problem (3A.8)

Suppose \(\lambda\) denotes Lebesgue measure on \(\mathbb{R}\). Give an example of a sequence \(f_1, f_2, \ldots\) of simple Borel measurable functions from \(\mathbb{R}\) to \([0, \infty)\) such that \(\lim_{k \to \infty} f_k(x) = 0\) for every \(x \in \mathbb{R}\) but \(\lim_{k \to \infty} \int f_k \, d\lambda = 1\).

Solution

For each \(k \in \mathbb{Z}^+\) define

\begin{equation*} f_k = \chi_{[k,\, k+1]}. \end{equation*}

Each \(f_k\) takes only the values \(0\) and \(1\), so each \(f_k\) is a simple function from \(\mathbb{R}\) to \([0, \infty)\), and each \(f_k\) is Borel measurable because \([k, k+1]\) is a Borel set.

Pointwise convergence to \(0\): fix \(x \in \mathbb{R}\) and choose \(N \in \mathbb{Z}^+\) with \(N > x\). If \(k \geq N\) then \(k > x\), so \(x \notin [k, k+1]\) and hence \(f_k(x) = 0\). Thus \(f_k(x) = 0\) for all \(k \geq N\), and in particular \(\lim_{k \to \infty} f_k(x) = 0\). Because \(x\) was arbitrary, the sequence converges pointwise to \(0\) on all of \(\mathbb{R}\).

Integrals: by 3.4 (the integral of a characteristic function is the measure of the set),

\begin{equation*} \int f_k \, d\lambda = \lambda\bigl([k, k+1]\bigr) = 1 \end{equation*}

for every \(k \in \mathbb{Z}^+\). Hence \(\lim_{k \to \infty} \int f_k \, d\lambda = 1\), while \(\int \bigl(\lim_{k \to \infty} f_k\bigr) d\lambda = \int 0 \, d\lambda = 0\).

This example shows that the hypothesis in the Monotone Convergence Theorem (3.11) that the sequence be increasing cannot simply be dropped: here the sequence is not monotone, and the limit of the integrals is not the integral of the limit.

Another example with the same behavior is \(f_k = \frac{1}{k}\chi_{[0,\,k]}\); then \(0 \leq f_k(x) \leq \frac{1}{k}\) for every \(x\), so \(f_k \to 0\) pointwise, while \(\int f_k \, d\lambda = \frac{1}{k} \cdot \lambda([0,k]) = 1\) for every \(k\). Note that in this second example the convergence to \(0\) is even uniform on \(\mathbb{R}\), which shows that uniform convergence alone does not allow the interchange of limits and integrals over a set of infinite measure.

Problem (3A.9)

Suppose \(\mu\) is a measure on a measurable space \((X, \mathcal{S})\) and \(f : X \to [0, \infty]\) is an \(\mathcal{S}\)-measurable function. Define \(\nu : \mathcal{S} \to [0, \infty]\) by

\begin{equation*} \nu(A) = \int \chi_A f \, d\mu \end{equation*}

for \(A \in \mathcal{S}\). Prove that \(\nu\) is a measure on \((X, \mathcal{S})\).

Solution

Throughout, \(\chi_A f\) denotes the function on \(X\) whose value at \(x\) is \(f(x)\) if \(x \in A\) and \(0\) if \(x \notin A\) (so no product of the form \(0 \cdot \infty\) arises; with the usual convention \(0 \cdot \infty = 0\) this is the pointwise product).

Step 1: \(\nu\) is well defined. Fix \(A \in \mathcal{S}\). The function \(\chi_A f\) maps \(X\) to \([0, \infty]\), and for \(a \in \mathbb{R}\) we have

\begin{equation*} (\chi_A f)^{-1}\bigl((a, \infty]\bigr) = \begin{cases} A \cap f^{-1}\bigl((a, \infty]\bigr) & \text{if } a \geq 0,\\ X & \text{if } a < 0, \end{cases} \end{equation*}

because \(\chi_A f\) vanishes off \(A\) and agrees with \(f\) on \(A\). Both of these sets are in \(\mathcal{S}\) (using that \(A \in \mathcal{S}\) and that \(f\) is \(\mathcal{S}\)-measurable), so \(\chi_A f\) is \(\mathcal{S}\)-measurable by 2.52. Hence \(\int \chi_A f \, d\mu\) is defined and lies in \([0, \infty]\), and \(\nu(A) \in [0, \infty]\).

Step 2: \(\nu(\emptyset) = 0\). We have \(\chi_\emptyset f = 0\), the function that is identically \(0\). For every \(\mathcal{S}\)-partition \(A_1, \ldots, A_m\) of \(X\) we get \(\mathcal{L}(0, P) = \sum_{j=1}^m \mu(A_j) \cdot 0 = 0\), so \(\int 0 \, d\mu = 0\). Thus \(\nu(\emptyset) = 0\).

Step 3: countable additivity. Suppose \(A_1, A_2, \ldots\) is a disjoint sequence of sets in \(\mathcal{S}\) and let \(A = \bigcup_{k=1}^\infty A_k\); note \(A \in \mathcal{S}\). For \(n \in \mathbb{Z}^+\) define

\begin{equation*} g_n = \sum_{k=1}^{n} \chi_{A_k} f = \chi_{A_1 \cup \cdots \cup A_n} f, \end{equation*}

where the second equality holds because the \(A_k\) are disjoint: at a point \(x\) lying in exactly one \(A_k\) with \(k \leq n\) the sum equals \(f(x)\), and at a point lying in none of \(A_1, \ldots, A_n\) every term is \(0\). By Step 1 (applied to the set \(A_1 \cup \cdots \cup A_n \in \mathcal{S}\)), each \(g_n\) is \(\mathcal{S}\)-measurable.

The sequence \(g_1, g_2, \ldots\) is increasing, since \(g_{n+1} - g_n = \chi_{A_{n+1}} f \geq 0\) pointwise (more precisely, \(g_{n+1}(x) = g_n(x)\) unless \(x \in A_{n+1}\), in which case \(g_n(x) = 0 \le f(x) = g_{n+1}(x)\)). Moreover

\begin{equation*} \lim_{n \to \infty} g_n(x) = (\chi_A f)(x) \qquad \text{for every } x \in X: \end{equation*}

if \(x \in A\) then \(x \in A_k\) for exactly one \(k\), and \(g_n(x) = f(x)\) for all \(n \geq k\); if \(x \notin A\) then \(g_n(x) = 0\) for all \(n\), and \((\chi_A f)(x) = 0\) as well.

The Monotone Convergence Theorem (3.11) now gives

\begin{equation*} \lim_{n \to \infty} \int g_n \, d\mu = \int \chi_A f \, d\mu = \nu(A). \end{equation*}

On the other hand, additivity of integration for nonnegative functions (3.16), applied \(n - 1\) times, gives

\begin{equation*} \int g_n \, d\mu = \int \Bigl( \sum_{k=1}^n \chi_{A_k} f \Bigr) d\mu = \sum_{k=1}^{n} \int \chi_{A_k} f \, d\mu = \sum_{k=1}^{n} \nu(A_k). \end{equation*}

Combining the two displays,

\begin{equation*} \nu(A) = \lim_{n \to \infty} \sum_{k=1}^{n} \nu(A_k) = \sum_{k=1}^{\infty} \nu(A_k), \end{equation*}

the last equality being the definition of a sum of a series in \([0, \infty]\) (the partial sums are increasing, so the limit exists in \([0, \infty]\)).

Steps 1 through 3 show that \(\nu\) is a function from \(\mathcal{S}\) to \([0, \infty]\) with \(\nu(\emptyset) = 0\) that is countably additive. Hence \(\nu\) is a measure on \((X, \mathcal{S})\), as required by 2.54.

Problem (3A.10)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(f_1, f_2, \ldots\) is a sequence of nonnegative \(\mathcal{S}\)-measurable functions. Define \(f : X \to [0, \infty]\) by \(f(x) = \sum_{k=1}^{\infty} f_k(x)\). Prove that

\begin{equation*} \int f \, d\mu = \sum_{k=1}^{\infty} \int f_k \, d\mu. \end{equation*}

Solution

For \(n \in \mathbb{Z}^+\) define the partial sum

\begin{equation*} g_n = \sum_{k=1}^{n} f_k . \end{equation*}

Measurability of \(g_n\). It suffices to check that the sum of two \(\mathcal{S}\)-measurable functions \(u, v : X \to [0, \infty]\) is \(\mathcal{S}\)-measurable, and then induct. For \(a < 0\) we have \((u + v)^{-1}((a, \infty]) = X \in \mathcal{S}\). For \(a \geq 0\) we claim

\begin{equation*} (u + v)^{-1}\bigl((a, \infty]\bigr) = \bigcup_{q \in \mathbb{Q}} \Bigl( u^{-1}\bigl((q, \infty]\bigr) \cap v^{-1}\bigl((a - q, \infty]\bigr) \Bigr). \end{equation*}

Indeed, if \(u(x) > q\) and \(v(x) > a - q\) then \(u(x) + v(x) > a\). Conversely suppose \(u(x) + v(x) > a\). If \(u(x)\) and \(v(x)\) are both finite, pick a rational \(q\) with \(a - v(x) < q < u(x)\); then \(u(x) > q\) and \(v(x) > a - q\). If \(u(x) = \infty\), pick a rational \(q > a\); then \(u(x) > q\) and \(v(x) \geq 0 > a - q\). If \(v(x) = \infty\), pick a rational \(q < 0\); then \(u(x) \geq 0 > q\) and \(v(x) = \infty > a - q\). The claimed identity holds, the right side is a countable union of sets in \(\mathcal{S}\), and so \(u + v\) is \(\mathcal{S}\)-measurable by 2.52. By induction each \(g_n\) is \(\mathcal{S}\)-measurable.

Monotonicity and limit. Because each \(f_k\) is nonnegative, \(g_{n+1}(x) - g_n(x) = f_{n+1}(x) \geq 0\), so \(0 \leq g_1 \leq g_2 \leq \cdots\). Moreover, by the definition of the sum of a series of nonnegative terms as the limit of its partial sums (a limit that always exists in \([0, \infty]\) because the partial sums increase),

\begin{equation*} \lim_{n \to \infty} g_n(x) = \sum_{k=1}^{\infty} f_k(x) = f(x) \qquad \text{for every } x \in X . \end{equation*}

In particular \(f\) is \(\mathcal{S}\)-measurable by 2.53 (it is the pointwise supremum of the \(g_n\)).

Applying the theorems. Additivity of integration for nonnegative functions (3.16), used \(n - 1\) times, gives

\begin{equation*} \int g_n \, d\mu = \sum_{k=1}^{n} \int f_k \, d\mu \end{equation*}

for each \(n \in \mathbb{Z}^+\). The Monotone Convergence Theorem (3.11) applies to the increasing sequence \(g_1 \le g_2 \le \cdots\) of nonnegative \(\mathcal{S}\)-measurable functions with pointwise limit \(f\), giving

\begin{equation*} \int f \, d\mu = \lim_{n \to \infty} \int g_n \, d\mu = \lim_{n \to \infty} \sum_{k=1}^{n} \int f_k \, d\mu = \sum_{k=1}^{\infty} \int f_k \, d\mu , \end{equation*}

where the last equality is again the definition of the sum of a series of terms in \([0, \infty]\). This is the desired equation; note that both sides may equal \(\infty\), and the identity holds in \([0, \infty]\) with no integrability hypothesis.

Problem (3A.11)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(f_1, f_2, \ldots\) are \(\mathcal{S}\)-measurable functions from \(X\) to \(\mathbb{R}\) such that \(\sum_{k=1}^{\infty} \int |f_k| \, d\mu < \infty\). Prove that there exists \(E \in \mathcal{S}\) such that \(\mu(X \setminus E) = 0\) and \(\lim_{k \to \infty} f_k(x) = 0\) for every \(x \in E\).

Solution

Each \(|f_k|\) is \(\mathcal{S}\)-measurable: for \(a < 0\) we have \(|f_k|^{-1}((a, \infty]) = X\), and for \(a \geq 0\)

\begin{equation*} |f_k|^{-1}\bigl((a, \infty]\bigr) = f_k^{-1}\bigl((a, \infty)\bigr) \cup f_k^{-1}\bigl((-\infty, -a)\bigr) \in \mathcal{S}, \end{equation*}

so 2.52 applies. (Alternatively, \(|f_k| = f_k^+ + f_k^-\) in the notation of 3.17.)

Define \(g : X \to [0, \infty]\) by

\begin{equation*} g(x) = \sum_{k=1}^{\infty} |f_k(x)| . \end{equation*}

The function \(g\) is \(\mathcal{S}\)-measurable: it is the pointwise supremum of the increasing sequence of partial sums \(\sum_{k=1}^{n} |f_k|\), each of which is \(\mathcal{S}\)-measurable (finite sums of \([0,\infty]\)-valued \(\mathcal{S}\)-measurable functions are \(\mathcal{S}\)-measurable, as shown in the solution to 3A.10), so 2.53 applies. Now the previous exercise (3A.10), applied to the nonnegative \(\mathcal{S}\)-measurable functions \(|f_1|, |f_2|, \ldots\), gives

\begin{equation*} \int g \, d\mu = \sum_{k=1}^{\infty} \int |f_k| \, d\mu < \infty , \end{equation*}

the finiteness being the hypothesis.

Now set

\begin{equation*} E = g^{-1}\bigl([0, \infty)\bigr) = \{ x \in X : g(x) < \infty \} . \end{equation*}

Since \(g\) is \(\mathcal{S}\)-measurable, \(E \in \mathcal{S}\), and \(X \setminus E = g^{-1}(\{\infty\}) \in \mathcal{S}\).

Claim: \(\mu(X \setminus E) = 0\). Fix \(t \in (0, \infty)\). For every \(x \in X\) we have

\begin{equation*} g(x) \geq t \, \chi_{X \setminus E}(x), \end{equation*}

because \(g(x) = \infty \geq t\) when \(x \in X \setminus E\), while the right side is \(0\) when \(x \in E\). Hence by 3.8 (integration is order preserving), 3.20 (homogeneity, or directly from the definition) and 3.4 (the integral of a characteristic function),

\begin{equation*} \int g \, d\mu \geq \int t \, \chi_{X \setminus E} \, d\mu = t \, \mu(X \setminus E). \end{equation*}

Therefore

\begin{equation*} \mu(X \setminus E) \leq \frac{1}{t} \int g \, d\mu \end{equation*}

for every \(t \in (0, \infty)\). Because \(\int g \, d\mu < \infty\), letting \(t \to \infty\) gives \(\mu(X \setminus E) = 0\).

Finally, suppose \(x \in E\). Then \(\sum_{k=1}^{\infty} |f_k(x)| = g(x) < \infty\), so this is a convergent series of nonnegative real numbers. Writing \(s_n = \sum_{k=1}^{n} |f_k(x)|\), the sequence \(s_1, s_2, \ldots\) converges to the real number \(g(x)\), so

\begin{equation*} |f_n(x)| = s_n - s_{n-1} \to g(x) - g(x) = 0 \end{equation*}

as \(n \to \infty\) (with \(s_0 = 0\)). Hence \(\lim_{k \to \infty} f_k(x) = 0\) for every \(x \in E\).

Thus \(E \in \mathcal{S}\) satisfies \(\mu(X \setminus E) = 0\) and \(f_k(x) \to 0\) for every \(x \in E\), completing the proof. In the language of Chapter 3, the conclusion says that \(f_k \to 0\) almost everywhere.

Problem (3A.12)

Show that there exists a Borel measurable function \(f : \mathbb{R} \to (0, \infty)\) such that \(\int \chi_I f \, d\lambda = \infty\) for every nonempty open interval \(I \subseteq \mathbb{R}\), where \(\lambda\) denotes Lebesgue measure on \(\mathbb{R}\).

Solution

The function must take values in \((0, \infty)\), so it must be finite at every point, yet fail to be integrable over every interval. We build it as a series of poles placed at the rationals, with coefficients decaying fast enough that the series converges off a set of measure \(0\); on that exceptional set we simply redefine the function to be \(1\).

Construction. Let \(r_1, r_2, \ldots\) be an enumeration of \(\mathbb{Q}\). For \(k \in \mathbb{Z}^+\) define \(h_k : \mathbb{R} \to (0, \infty]\) by

\begin{equation*} h_k(x) = \frac{1}{|x - r_k|} \quad \text{if } x \neq r_k, \qquad h_k(r_k) = \infty . \end{equation*}

Each \(h_k\) is Borel measurable: if \(a > 0\) then \(h_k^{-1}((a, \infty]) = (r_k - \frac{1}{a},\, r_k + \frac{1}{a})\), which is open, and if \(a \leq 0\) then \(h_k^{-1}((a, \infty]) = \mathbb{R}\); now apply 2.52. Define \(g : \mathbb{R} \to (0, \infty]\) by

\begin{equation*} g(x) = \sum_{k=1}^{\infty} 8^{-k} h_k(x). \end{equation*}

Each partial sum is Borel measurable and \(g\) is their pointwise supremum, so \(g\) is Borel measurable by 2.53. Every term is strictly positive, so \(g(x) > 0\) for all \(x\).

The set where \(g\) is infinite is null. Let

\begin{equation*} B = \bigcap_{n=1}^{\infty} \; \bigcup_{k=n}^{\infty} \bigl( r_k - 4^{-k},\, r_k + 4^{-k} \bigr), \end{equation*}

a Borel set, consisting of those \(x\) with \(|x - r_k| < 4^{-k}\) for infinitely many \(k\). For each \(n\), monotonicity and countable subadditivity of \(\lambda\) (2.57 and 2.58) give

\begin{equation*} \lambda(B) \leq \sum_{k=n}^{\infty} 2 \cdot 4^{-k} = \frac{8}{3} \cdot 4^{-n} . \end{equation*}

Letting \(n \to \infty\) yields \(\lambda(B) = 0\).

Now suppose \(x \notin B \cup \mathbb{Q}\). Since \(x \notin B\) there is \(N \in \mathbb{Z}^+\) with \(|x - r_k| \geq 4^{-k}\) for all \(k \geq N\), whence

\begin{equation*} 8^{-k} h_k(x) = \frac{8^{-k}}{|x - r_k|} \leq 8^{-k} \cdot 4^{k} = 2^{-k} \qquad (k \geq N). \end{equation*}

Since \(x \notin \mathbb{Q}\), each of the finitely many terms with \(k < N\) is a finite real number. Therefore

\begin{equation*} g(x) \leq \sum_{k=1}^{N-1} 8^{-k} h_k(x) + \sum_{k=N}^{\infty} 2^{-k} < \infty . \end{equation*}

Consequently the Borel set \(S = g^{-1}(\{\infty\})\) satisfies \(S \subseteq B \cup \mathbb{Q}\), and hence

\begin{equation*} \lambda(S) \leq \lambda(B) + \lambda(\mathbb{Q}) = 0 . \end{equation*}

Definition of \(f\). Define \(f : \mathbb{R} \to (0, \infty)\) by

\begin{equation*} f(x) = \begin{cases} g(x) & \text{if } x \notin S, \\ 1 & \text{if } x \in S. \end{cases} \end{equation*}

For \(x \notin S\) we have \(0 < g(x) < \infty\), so indeed \(f\) takes values in \((0, \infty)\). For \(a \in \mathbb{R}\),

\begin{equation*} f^{-1}\bigl((a, \infty)\bigr) = \Bigl( g^{-1}\bigl((a, \infty]\bigr) \setminus S \Bigr) \cup \begin{cases} S & \text{if } a < 1, \\ \emptyset & \text{if } a \geq 1, \end{cases} \end{equation*}

which is a Borel set; so \(f\) is Borel measurable by 2.52.

Verification. Fix a nonempty open interval \(I \subseteq \mathbb{R}\).

Step 1: \(\int \chi_I h_j \, d\lambda = \infty\) whenever \(r_j \in I\). Choose \(\delta > 0\) with \((r_j - \delta, r_j + \delta) \subseteq I\). For \(m = 0, 1, 2, \ldots\) put

\begin{equation*} A_m = \bigl( r_j + \delta 2^{-m-2},\; r_j + \delta 2^{-m-1} \bigr] . \end{equation*}

These are disjoint Borel subsets of \((r_j, r_j + \delta/2] \subseteq I\), with \(\lambda(A_m) = \delta 2^{-m-2}\). For \(x \in A_m\) we have \(0 < x - r_j \leq \delta 2^{-m-1}\), so \(h_j(x) \geq 2^{m+1}/\delta\). Hence for each \(M\),

\begin{equation*} \sum_{m=0}^{M} \frac{2^{m+1}}{\delta} \, \chi_{A_m}(x) \leq \chi_I(x) h_j(x) \qquad \text{for all } x \in \mathbb{R}, \end{equation*}

so 3.9 gives

\begin{equation*} \int \chi_I h_j \, d\lambda \geq \sum_{m=0}^{M} \frac{2^{m+1}}{\delta} \cdot \delta 2^{-m-2} = \frac{M+1}{2}. \end{equation*}

Letting \(M \to \infty\) proves Step 1.

Step 2: \(\int \chi_I g \, d\lambda = \infty\). Since \(\mathbb{Q}\) is dense and \(I\) is a nonempty open interval, there is \(j\) with \(r_j \in I\). The nonnegative Borel functions \(8^{-k}\chi_I h_k\) sum pointwise to \(\chi_I g\), so by 3A.10 and homogeneity of integration (3.20),

\begin{equation*} \int \chi_I g \, d\lambda = \sum_{k=1}^{\infty} 8^{-k} \int \chi_I h_k \, d\lambda \geq 8^{-j} \int \chi_I h_j \, d\lambda = \infty . \end{equation*}

Step 3: \(\int \chi_{I \cap S} \, g \, d\lambda = 0\). The function \(\chi_{I \cap S} \, g\) equals \(\infty\) on \(I \cap S\) and \(0\) elsewhere. Suppose \(A_1, \ldots, A_m\) are disjoint Borel sets and \(c_1, \ldots, c_m \in [0, \infty)\) satisfy \(\sum_{j=1}^m c_j \chi_{A_j} \leq \chi_{I \cap S} \, g\) pointwise. If \(c_j > 0\) and \(x \in A_j\), then the left side is positive at \(x\), forcing \(x \in I \cap S\); thus \(A_j \subseteq S\) and \(\lambda(A_j) = 0\). Therefore \(\sum_{j=1}^m c_j \lambda(A_j) = 0\) for every such choice, and 3.9 gives \(\int \chi_{I \cap S} \, g \, d\lambda = 0\).

Step 4: conclusion. Pointwise \(\chi_I g = \chi_{I \cap S} \, g + \chi_{I \setminus S} \, g\), so additivity of integration (3.16) together with Steps 2 and 3 gives

\begin{equation*} \infty = \int \chi_I g \, d\lambda = 0 + \int \chi_{I \setminus S} \, g \, d\lambda , \end{equation*}

so \(\int \chi_{I \setminus S} \, g \, d\lambda = \infty\). Finally \(\chi_I f \geq \chi_{I \setminus S} \, g\) pointwise, since the two agree on \(I \setminus S\) and the right side is \(0\) off \(I \setminus S\). Hence by 3.8,

\begin{equation*} \int \chi_I f \, d\lambda = \infty . \end{equation*}

Because \(I\) was an arbitrary nonempty open interval, the Borel measurable function \(f : \mathbb{R} \to (0, \infty)\) has the required property.

Problem (3A.13)

Give an example to show that the Monotone Convergence Theorem (3.11) can fail if the hypothesis that \(f_1, f_2, \ldots\) are nonnegative functions is dropped.

Solution

Let \(\lambda\) denote Lebesgue measure on \(\mathbb{R}\), with \(\mathcal{S}\) the \(\sigma\)-algebra of Borel subsets of \(\mathbb{R}\). For \(k \in \mathbb{Z}^+\) define \(f_k : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} f_k = -\chi_{[k, \infty)}, \qquad \text{that is,} \qquad f_k(x) = \begin{cases} -1 & \text{if } x \geq k, \\ 0 & \text{if } x < k. \end{cases} \end{equation*}

Each \(f_k\) is a Borel measurable simple function.

The sequence is increasing. If \(x < k\) then \(f_k(x) = 0 = f_{k+1}(x)\). If \(k \leq x < k+1\) then \(f_k(x) = -1 < 0 = f_{k+1}(x)\). If \(x \geq k+1\) then \(f_k(x) = -1 = f_{k+1}(x)\). In every case \(f_k(x) \leq f_{k+1}(x)\), so \(f_1 \leq f_2 \leq \cdots\). What fails is only the hypothesis \(f_k \geq 0\).

The pointwise limit is \(0\). Fix \(x \in \mathbb{R}\) and take \(N \in \mathbb{Z}^+\) with \(N > x\). Then \(f_k(x) = 0\) for all \(k \geq N\), so

\begin{equation*} f(x) := \lim_{k \to \infty} f_k(x) = 0 \qquad \text{for every } x \in \mathbb{R}. \end{equation*}

Hence \(\int f \, d\lambda = \int 0 \, d\lambda = 0\).

The integrals. In the notation of 3.17, \(f_k^{\,+} = 0\) and \(f_k^{\,-} = \chi_{[k,\infty)}\). Thus \(\int f_k^{\,+} d\lambda = 0\) is finite, so \(\int f_k \, d\lambda\) is defined by 3.18, and by 3.4

\begin{equation*} \int f_k \, d\lambda = \int f_k^{\,+} \, d\lambda - \int f_k^{\,-} \, d\lambda = 0 - \lambda\bigl([k, \infty)\bigr) = -\infty . \end{equation*}

Therefore

\begin{equation*} \lim_{k \to \infty} \int f_k \, d\lambda = -\infty \neq 0 = \int f \, d\lambda , \end{equation*}

and the conclusion of the Monotone Convergence Theorem fails for this increasing sequence of measurable functions that are not nonnegative.

Remark. The failure is caused entirely by the infinite negative mass at the bottom of the sequence, and no example can avoid that. Indeed, suppose \(f_1 \leq f_2 \leq \cdots\) are \(\mathcal{S}\)-measurable real-valued functions with \(\int |f_1| \, d\mu < \infty\). Then \(f_k - f_1\) is a nonnegative increasing sequence converging pointwise to \(f - f_1\), so 3.11 gives \(\int (f_k - f_1) \, d\mu \to \int (f - f_1) \, d\mu\); adding the finite number \(\int f_1 \, d\mu\) to both sides recovers \(\int f_k \, d\mu \to \int f \, d\mu\). (This is the increasing half of Exercise 20 in this section.) In particular, in any counterexample \(\int |f_1| \, d\mu\) must be infinite, as it is above: \(\int |f_1| \, d\lambda = \lambda([1, \infty)) = \infty\). So the nonnegativity hypothesis in 3.11 cannot simply be deleted, but it can be weakened to \(\int |f_1| \, d\mu < \infty\).

Problem (3A.14)

Give an example to show that the Monotone Convergence Theorem can fail if the hypothesis of an increasing sequence of functions is replaced by a hypothesis of a decreasing sequence of functions.

[This exercise shows that the Monotone Convergence Theorem should be called the Increasing Convergence Theorem. However, see Exercise 20.]

Solution

Let \(\lambda\) denote Lebesgue measure on \(\mathbb{R}\), with the Borel \(\sigma\)-algebra. For \(k \in \mathbb{Z}^+\) define

\begin{equation*} f_k = \chi_{[k, \infty)}, \qquad \text{that is,} \qquad f_k(x) = \begin{cases} 1 & \text{if } x \geq k, \\ 0 & \text{if } x < k. \end{cases} \end{equation*}

Each \(f_k\) is a nonnegative Borel measurable simple function, so every hypothesis of 3.11 holds except that the sequence decreases rather than increases.

The sequence is decreasing. Since \([k+1, \infty) \subseteq [k, \infty)\), we have \(\chi_{[k+1,\infty)} \leq \chi_{[k,\infty)}\) pointwise, that is \(f_1 \geq f_2 \geq \cdots \geq 0\).

The pointwise limit is \(0\). Fix \(x \in \mathbb{R}\) and choose \(N \in \mathbb{Z}^+\) with \(N > x\). Then \(x \notin [k, \infty)\) for every \(k \geq N\), so \(f_k(x) = 0\) for all \(k \geq N\). Hence

\begin{equation*} f(x) := \lim_{k \to \infty} f_k(x) = 0 \qquad \text{for every } x \in \mathbb{R}, \end{equation*}

and therefore \(\int f \, d\lambda = 0\).

The integrals. By 3.4,

\begin{equation*} \int f_k \, d\lambda = \lambda\bigl([k, \infty)\bigr) = \infty \end{equation*}

for every \(k \in \mathbb{Z}^+\). Thus

\begin{equation*} \lim_{k \to \infty} \int f_k \, d\lambda = \infty \neq 0 = \int f \, d\lambda , \end{equation*}

so the conclusion of the Monotone Convergence Theorem fails for this decreasing sequence.

Remark. The failure is driven entirely by the infinite mass of \(f_1\), and it must be. If \(f_1 \geq f_2 \geq \cdots \geq 0\) are \(\mathcal{S}\)-measurable with \(f_1\) real-valued and \(\int f_1 \, d\mu < \infty\), then \(f_1 - f_1, f_1 - f_2, \ldots\) is an increasing sequence of nonnegative \(\mathcal{S}\)-measurable functions with pointwise limit \(f_1 - f\), so 3.11 gives \(\int (f_1 - f_k) \, d\mu \to \int (f_1 - f) \, d\mu\); subtracting this from the finite number \(\int f_1 \, d\mu\) yields \(\int f_k \, d\mu \to \int f \, d\mu\). This is the decreasing half of Exercise 20 in this section, and its hypothesis \(\int |f_1| \, d\mu < \infty\) fails in the example above, where \(\int f_1 \, d\lambda = \infty\).

In particular, no counterexample can have \(f_1\) bounded on a space of finite measure: if \(\mu(X) < \infty\) and \(c = \sup_X f_1 < \infty\), then \(f_1 \leq c \, \chi_X\) pointwise, so 3.8 and 3.4 give \(\int f_1 \, d\mu \leq c \, \mu(X) < \infty\). But finite measure alone is not enough. Let \(X = (0, 1]\) with the Borel \(\sigma\)-algebra and Lebesgue measure, and define

\begin{equation*} f_k(x) = \begin{cases} 1/x & \text{if } 0 < x \leq 1/k, \\ 0 & \text{if } 1/k < x \leq 1. \end{cases} \end{equation*}

Each \(f_k\) is Borel measurable, and \(f_1 \geq f_2 \geq \cdots \geq 0\) because the sets \((0, 1/k]\) decrease while the formula \(1/x\) does not depend on \(k\). For each \(x \in (0,1]\) we have \(f_k(x) = 0\) once \(1/k < x\), so \(f_k \to 0\) pointwise. Yet \(\int f_k \, d\lambda = \infty\) for every \(k\): the sets \(A_m = (2^{-m-1}/k, \; 2^{-m}/k]\) for \(m = 0, 1, 2, \ldots\) are disjoint subsets of \((0, 1/k]\) with \(\lambda(A_m) = 2^{-m-1}/k\), and \(f_k \geq k \, 2^{m}\) on \(A_m\), so \(\sum_{m=0}^{M} k \, 2^{m} \chi_{A_m} \leq f_k\) pointwise and 3.9 gives \(\int f_k \, d\lambda \geq \sum_{m=0}^{M} k \, 2^{m} \cdot 2^{-m-1}/k = (M+1)/2\) for every \(M\). Here \(\lambda\bigl((0,1]\bigr) = 1 < \infty\), but \(f_1\) is unbounded and \(\int f_1 \, d\lambda = \infty\), so Exercise 20 again does not apply.

Problem (3A.15)

Suppose \(\lambda\) is Lebesgue measure on \(\mathbb{R}\) and \(f : \mathbb{R} \to [-\infty, \infty]\) is a Borel measurable function such that \(\int f \, d\lambda\) is defined.

(a) For \(t \in \mathbb{R}\), define \(f_t : \mathbb{R} \to [-\infty, \infty]\) by \(f_t(x) = f(x - t)\). Prove that \(\int f_t \, d\lambda = \int f \, d\lambda\) for all \(t \in \mathbb{R}\).

(b) For \(t \in \mathbb{R}\), define \(f_t : \mathbb{R} \to [-\infty, \infty]\) by \(f_t(x) = f(tx)\). Prove that \(\int f_t \, d\lambda = \frac{1}{|t|} \int f \, d\lambda\) for all \(t \in \mathbb{R} \setminus \{0\}\).

Solution

Throughout, \(\mathcal{B}\) denotes the \(\sigma\)-algebra of Borel subsets of \(\mathbb{R}\), and \(\lambda\) is Lebesgue measure on \(\mathcal{B}\); thus \(\lambda(A) = |A|\) (outer measure) for every \(A \in \mathcal{B}\), by 2.68 and 2.69.

Preliminaries on the two changes of variable.

Fix \(t \in \mathbb{R}\) and let \(\tau : \mathbb{R} \to \mathbb{R}\) be \(\tau(x) = x - t\). For \(A \subseteq \mathbb{R}\) we have \(\tau^{-1}(A) = t + A\). Because \(\tau\) is continuous it is Borel measurable (2.41), so

\begin{equation*} A \in \mathcal{B} \implies t + A \in \mathcal{B}. \end{equation*}

Moreover \(A \mapsto t + A\) is a bijection of \(\mathcal{B}\) onto \(\mathcal{B}\) (its inverse is \(A \mapsto -t + A\)), it preserves disjointness, and it satisfies \((t + A_1) \cup \cdots \cup (t + A_m) = t + (A_1 \cup \cdots \cup A_m)\). Finally \(\lambda(t + A) = \lambda(A)\) for all \(A \in \mathcal{B}\), by 2.7.

Now fix \(t \in \mathbb{R} \setminus \{0\}\) and let \(\sigma(x) = tx\), so that \(\sigma^{-1}(A) = \frac{1}{t} A\), where \(cA := \{ca : a \in A\}\). Again \(\sigma\) is continuous, hence Borel measurable, so \(A \in \mathcal{B}\) implies \(\frac{1}{t}A \in \mathcal{B}\). The map \(A \mapsto \frac{1}{t}A\) is a bijection of \(\mathcal{B}\) onto \(\mathcal{B}\) (inverse \(A \mapsto tA\)), preserves disjointness and unions, and by Exercise 2 in Section 2A we have \(|cA| = |c| \, |A|\) for all \(c \in \mathbb{R}\) and \(A \subseteq \mathbb{R}\); hence

\begin{equation*} \lambda\Bigl(\tfrac{1}{t}A\Bigr) = \frac{1}{|t|}\,\lambda(A) \qquad \text{for all } A \in \mathcal{B}. \end{equation*}

Measurability of \(f_t\). In case (a), for every Borel subset \(B\) of \([-\infty,\infty]\),

\begin{equation*} f_t^{-1}(B) = \{x : f(x - t) \in B\} = t + f^{-1}(B), \end{equation*}

which is a Borel subset of \(\mathbb{R}\) by the above. Thus \(f_t\) is Borel measurable. In case (b), \(f_t^{-1}(B) = \{x : f(tx) \in B\} = \frac{1}{t} f^{-1}(B)\), again Borel. So in both cases \(f_t\) is Borel measurable.

(a) The nonnegative case. Assume first \(f \ge 0\). Let \(P\) be a Borel partition \(A_1, \dots, A_m\) of \(\mathbb{R}\), and let \(P_t\) denote \(t + A_1, \dots, t + A_m\), which by the preliminaries is again a Borel partition of \(\mathbb{R}\). Because \(x \in t + A_j\) if and only if \(x - t \in A_j\),

\begin{equation*} \inf_{t + A_j} f_t = \inf\{f(x - t) : x \in t + A_j\} = \inf\{f(y) : y \in A_j\} = \inf_{A_j} f . \end{equation*}

Combining this with \(\lambda(t + A_j) = \lambda(A_j)\) gives

\begin{equation*} L(f_t, P_t) = \sum_{j=1}^{m} \lambda(t + A_j) \inf_{t+A_j} f_t = \sum_{j=1}^{m} \lambda(A_j) \inf_{A_j} f = L(f, P). \end{equation*}

Since \(P \mapsto P_t\) is a bijection from the set of Borel partitions of \(\mathbb{R}\) onto itself, the two sets of lower Lebesgue sums \(\{L(f,P) : P\}\) and \(\{L(f_t, Q) : Q\}\) coincide. Taking suprema and using the definition 3.3 of the integral,

\begin{equation*} \int f_t \, d\lambda = \int f \, d\lambda . \end{equation*}

(a) The general case. For arbitrary \(f : \mathbb{R} \to [-\infty,\infty]\) note the pointwise identities

\begin{equation*} (f_t)^+ (x) = \max\{f(x-t), 0\} = (f^+)_t(x), \qquad (f_t)^-(x) = (f^-)_t(x). \end{equation*}

By the nonnegative case applied to \(f^+\) and to \(f^-\),

\begin{equation*} \int (f_t)^+ \, d\lambda = \int f^+ \, d\lambda \qquad \text{and} \qquad \int (f_t)^- \, d\lambda = \int f^- \, d\lambda . \end{equation*}

Because \(\int f \, d\lambda\) is defined, at least one of \(\int f^+ \, d\lambda\), \(\int f^- \, d\lambda\) is finite; hence at least one of \(\int (f_t)^+ \, d\lambda\), \(\int (f_t)^- \, d\lambda\) is finite, so \(\int f_t \, d\lambda\) is defined, and by 3.18

\begin{equation*} \int f_t \, d\lambda = \int (f_t)^+ d\lambda - \int (f_t)^- d\lambda = \int f^+ d\lambda - \int f^- d\lambda = \int f \, d\lambda . \end{equation*}

(b) The nonnegative case. Fix \(t \ne 0\) and assume \(f \ge 0\). Let \(P\) be a Borel partition \(A_1, \dots, A_m\) of \(\mathbb{R}\) and let \(P^{(t)}\) denote \(\frac{1}{t}A_1, \dots, \frac{1}{t}A_m\), a Borel partition of \(\mathbb{R}\) by the preliminaries. Since \(x \in \frac{1}{t}A_j\) if and only if \(tx \in A_j\),

\begin{equation*} \inf_{\frac{1}{t}A_j} f_t = \inf\{f(tx) : x \in \tfrac{1}{t}A_j\} = \inf\{f(y) : y \in A_j\} = \inf_{A_j} f . \end{equation*}

Therefore

\begin{equation*} L\bigl(f_t, P^{(t)}\bigr) = \sum_{j=1}^{m} \lambda\Bigl(\tfrac{1}{t}A_j\Bigr) \inf_{A_j} f = \frac{1}{|t|} \sum_{j=1}^{m} \lambda(A_j) \inf_{A_j} f = \frac{1}{|t|} L(f, P). \end{equation*}

As \(P \mapsto P^{(t)}\) is a bijection of the set of Borel partitions of \(\mathbb{R}\) onto itself, the set of lower Lebesgue sums for \(f_t\) is exactly \(\frac{1}{|t|}\) times the set of lower Lebesgue sums for \(f\). Since multiplication by the positive constant \(\frac{1}{|t|}\) commutes with taking suprema in \([0,\infty]\), definition 3.3 gives

\begin{equation*} \int f_t \, d\lambda = \frac{1}{|t|} \int f \, d\lambda . \end{equation*}

(b) The general case. Again \((f_t)^{\pm} = (f^{\pm})_t\) pointwise, so by the nonnegative case

\begin{equation*} \int (f_t)^{\pm} \, d\lambda = \frac{1}{|t|} \int f^{\pm} \, d\lambda . \end{equation*}

Since \(\frac{1}{|t|} \in (0,\infty)\), one of these is finite exactly when the corresponding integral of \(f^{\pm}\) is finite; hence \(\int f_t \, d\lambda\) is defined and

\begin{equation*} \int f_t \, d\lambda = \frac{1}{|t|} \int f^+ d\lambda - \frac{1}{|t|} \int f^- d\lambda = \frac{1}{|t|} \int f \, d\lambda , \end{equation*}

the last step being valid in \([-\infty,\infty]\) because at most one of the two terms is infinite.

Problem (3A.16)

Suppose \(\mathcal{S}\) and \(\mathcal{T}\) are \(\sigma\)-algebras on a set \(X\) and \(\mathcal{S} \subseteq \mathcal{T}\). Suppose \(\mu_1\) is a measure on \((X, \mathcal{S})\), \(\mu_2\) is a measure on \((X, \mathcal{T})\), and \(\mu_1(E) = \mu_2(E)\) for all \(E \in \mathcal{S}\). Prove that if \(f : X \to [0, \infty]\) is \(\mathcal{S}\)-measurable, then \(\int f \, d\mu_1 = \int f \, d\mu_2\).

Solution

First note that every \(\mathcal{S}\)-measurable function is \(\mathcal{T}\)-measurable: if \(g\) is \(\mathcal{S}\)-measurable and \(B\) is a Borel subset of \([-\infty,\infty]\), then \(g^{-1}(B) \in \mathcal{S} \subseteq \mathcal{T}\). In particular \(f\) is \(\mathcal{T}\)-measurable, so \(\int f \, d\mu_2\) makes sense.

One inequality is immediate. Every \(\mathcal{S}\)-partition \(A_1, \dots, A_m\) of \(X\) is also a \(\mathcal{T}\)-partition of \(X\), and since \(\mu_1(A_j) = \mu_2(A_j)\) for each \(j\) the two lower Lebesgue sums agree:

\begin{equation*} \sum_{j=1}^{m} \mu_1(A_j) \inf_{A_j} f = \sum_{j=1}^{m} \mu_2(A_j) \inf_{A_j} f . \end{equation*}

Hence the supremum defining \(\int f \, d\mu_2\) is taken over a collection of numbers containing all the lower Lebesgue sums that define \(\int f \, d\mu_1\), so by 3.3

\begin{equation*} \int f \, d\mu_1 \le \int f \, d\mu_2 . \end{equation*}

The reverse inequality is the real content, and it is where the \(\mathcal{S}\)-measurability of \(f\) (as opposed to mere \(\mathcal{T}\)-measurability) is used. We prove the equality outright by approximation.

Step 1: the result for simple \(\mathcal{S}\)-measurable functions. Suppose \(h : X \to [0,\infty)\) is a simple \(\mathcal{S}\)-measurable function. Write \(h\) in its standard representation

\begin{equation*} h = \sum_{k=1}^{n} c_k \chi_{E_k}, \end{equation*}

where \(c_1, \dots, c_n \in [0,\infty)\) are the distinct values of \(h\) and \(E_k = h^{-1}(\{c_k\})\); thus \(E_1, \dots, E_n\) are disjoint sets in \(\mathcal{S}\) whose union is \(X\). Because \(\mathcal{S} \subseteq \mathcal{T}\), these sets also lie in \(\mathcal{T}\), so 3.7 applies in each of the two measure spaces and gives

\begin{equation*} \int h \, d\mu_1 = \sum_{k=1}^{n} c_k \, \mu_1(E_k) = \sum_{k=1}^{n} c_k \, \mu_2(E_k) = \int h \, d\mu_2 , \end{equation*}

the middle equality holding because each \(E_k \in \mathcal{S}\), where \(\mu_1\) and \(\mu_2\) agree.

Step 2: an increasing sequence of nonnegative simple \(\mathcal{S}\)-measurable functions converging to \(f\). By 2.89, applied to the measurable space \((X, \mathcal{S})\) and to \(f\), there is a sequence \(u_1, u_2, \dots\) of simple \(\mathcal{S}\)-measurable functions from \(X\) to \(\mathbb{R}\) such that \(|u_k(x)| \le |u_{k+1}(x)| \le |f(x)|\) for all \(k \in \mathbb{Z}^+\) and all \(x \in X\), and \(\lim_{k \to \infty} u_k(x) = f(x)\) for every \(x \in X\). Define

\begin{equation*} f_k = \max\{u_1^+, u_2^+, \dots, u_k^+\} \qquad (k \in \mathbb{Z}^+), \end{equation*}

where \(u^+ = \max\{u, 0\}\). Each \(f_k\) takes only finitely many values (each value is a value of some \(u_j^+\)), and for \(a \in \mathbb{R}\) we have \(f_k^{-1}\bigl((a,\infty]\bigr) = \bigcup_{j=1}^{k} (u_j^+)^{-1}\bigl((a,\infty]\bigr) \in \mathcal{S}\), so by 2.52 each \(f_k\) is a simple nonnegative \(\mathcal{S}\)-measurable function. By construction \(f_1 \le f_2 \le \cdots\). Since \(f \ge 0\) we have \(|f| = f\), so \(u_j^+ \le |u_j| \le f\) for every \(j\), whence \(f_k \le f\); and \(f_k \ge u_k^+ \to f^+ = f\) pointwise. Therefore

\begin{equation*} 0 \le f_1 \le f_2 \le \cdots \quad \text{and} \quad \lim_{k \to \infty} f_k(x) = f(x) \ \text{ for every } x \in X . \end{equation*}

(This is the same use of 2.89 that the book makes in its proof of 3.16.)

Each \(f_k\) is also a simple \(\mathcal{T}\)-measurable function, and \(f\) is \(\mathcal{T}\)-measurable. Hence the Monotone Convergence Theorem 3.11 applies both in the measure space \((X, \mathcal{S}, \mu_1)\) and in the measure space \((X, \mathcal{T}, \mu_2)\). Using Step 1 in the middle,

\begin{equation*} \int f \, d\mu_1 = \lim_{k \to \infty} \int f_k \, d\mu_1 = \lim_{k \to \infty} \int f_k \, d\mu_2 = \int f \, d\mu_2 , \end{equation*}

which completes the proof.

Problem (3A.17)

Suppose that \((X, \mathcal{S}, \mu)\) is a measure space and \(f_1, f_2, \dots\) is a sequence of nonnegative \(\mathcal{S}\)-measurable functions on \(X\). Define a function \(f : X \to [0, \infty]\) by

\begin{equation*} f(x) = \liminf_{k \to \infty} f_k(x). \end{equation*}

(a) Show that \(f\) is an \(\mathcal{S}\)-measurable function.

(b) Prove that

\begin{equation*} \int f \, d\mu \le \liminf_{k \to \infty} \int f_k \, d\mu . \end{equation*}

(c) Give an example showing that the inequality in (b) can be a strict inequality even when \(\mu(X) < \infty\) and the family of functions \(\{f_k\}_{k \in \mathbb{Z}^+}\) is uniformly bounded.

[The result in (b) is called Fatou’s Lemma. Some textbooks prove Fatou’s Lemma and then use it to prove the Monotone Convergence Theorem. Here we are taking the reverse approach—you should be able to use the Monotone Convergence Theorem to give a clean proof of Fatou’s Lemma.]

Solution

For each \(k \in \mathbb{Z}^+\) define \(g_k : X \to [0, \infty]\) by

\begin{equation*} g_k(x) = \inf\{f_j(x) : j \ge k\} = \inf\{f_k(x), f_{k+1}(x), \dots\}. \end{equation*}

(a) Each \(g_k\) is the pointwise infimum of the sequence of \(\mathcal{S}\)-measurable functions \(f_k, f_{k+1}, \dots\), hence \(g_k\) is \(\mathcal{S}\)-measurable by 2.53. For each \(x \in X\) the sequence \(g_1(x) \le g_2(x) \le \cdots\) is increasing (the infimum is taken over a smaller set as \(k\) grows), so

\begin{equation*} f(x) = \liminf_{k \to \infty} f_k(x) = \lim_{k \to \infty} g_k(x) = \sup\{g_k(x) : k \in \mathbb{Z}^+\} . \end{equation*}

Thus \(f\) is the pointwise supremum of the sequence \(g_1, g_2, \dots\) of \(\mathcal{S}\)-measurable functions, so \(f\) is \(\mathcal{S}\)-measurable by 2.53.

(b) By part (a), \(0 \le g_1 \le g_2 \le \cdots\) is an increasing sequence of nonnegative \(\mathcal{S}\)-measurable functions with \(\lim_{k \to \infty} g_k(x) = f(x)\) for every \(x \in X\). Hence the Monotone Convergence Theorem 3.11 gives

\begin{equation*} \int f \, d\mu = \lim_{k \to \infty} \int g_k \, d\mu . \end{equation*}

Fix \(k \in \mathbb{Z}^+\). For every \(j \ge k\) we have \(g_k(x) \le f_j(x)\) for all \(x \in X\), so 3.8 gives \(\int g_k \, d\mu \le \int f_j \, d\mu\) for every \(j \ge k\). Taking the infimum over \(j \ge k\),

\begin{equation*} \int g_k \, d\mu \le \inf\Bigl\{ \int f_j \, d\mu : j \ge k \Bigr\} . \end{equation*}

Now let \(k \to \infty\). The left side tends to \(\int f \, d\mu\), and by the definition of \(\liminf\) the right side tends to \(\liminf_{k \to \infty} \int f_k \, d\mu\). Since the inequality is preserved in the limit,

\begin{equation*} \int f \, d\mu \le \liminf_{k \to \infty} \int f_k \, d\mu , \end{equation*}

as desired.

(c) Let \(X = [0,1]\), let \(\mathcal{S}\) be the \(\sigma\)-algebra of Borel subsets of \([0,1]\), and let \(\mu = \lambda\) be Lebesgue measure on \(\mathcal{S}\), so \(\mu(X) = 1 < \infty\). Define

\begin{equation*} f_k = \begin{cases} \chi_{[0, 1/2]} & \text{if } k \text{ is odd,} \\ \chi_{(1/2, \, 1]} & \text{if } k \text{ is even.} \end{cases} \end{equation*}

Each \(f_k\) is a nonnegative Borel measurable function with \(0 \le f_k \le 1\), so the family \(\{f_k\}_{k \in \mathbb{Z}^+}\) is uniformly bounded.

For each \(x \in [0,1]\), the sequence \(f_1(x), f_2(x), \dots\) takes the value \(0\) for infinitely many \(k\) (namely for all even \(k\) if \(x \in [0, 1/2]\), and for all odd \(k\) if \(x \in (1/2, 1]\)). Hence \(\inf\{f_j(x) : j \ge k\} = 0\) for every \(k\), and therefore

\begin{equation*} f(x) = \liminf_{k \to \infty} f_k(x) = 0 \qquad \text{for every } x \in [0,1]. \end{equation*}

Consequently \(\int f \, d\lambda = 0\). On the other hand, by 3.4,

\begin{equation*} \int f_k \, d\lambda = \lambda\bigl([0, \tfrac{1}{2}]\bigr) = \tfrac{1}{2} \quad \text{for odd } k, \qquad \int f_k \, d\lambda = \lambda\bigl((\tfrac{1}{2}, 1]\bigr) = \tfrac{1}{2} \quad \text{for even } k . \end{equation*}

Thus \(\liminf_{k \to \infty} \int f_k \, d\lambda = \frac{1}{2}\), and

\begin{equation*} 0 = \int f \, d\lambda < \liminf_{k \to \infty} \int f_k \, d\lambda = \tfrac{1}{2}, \end{equation*}

so the inequality in (b) is strict here.

Problem (3A.18)

Give an example of a sequence \(x_1, x_2, \dots\) of real numbers such that

\begin{equation*} \lim_{n \to \infty} \sum_{k=1}^{n} x_k \ \text{ exists in } \mathbb{R}, \end{equation*}

but \(\int x \, d\mu\) is not defined, where \(\mu\) is counting measure on \(\mathbb{Z}^+\) and \(x\) is the function from \(\mathbb{Z}^+\) to \(\mathbb{R}\) defined by \(x(k) = x_k\).

Solution

Take

\begin{equation*} x_k = \frac{(-1)^k}{k} \qquad (k \in \mathbb{Z}^+), \end{equation*}

so that \(x_1 = -1\), \(x_2 = \frac{1}{2}\), \(x_3 = -\frac{1}{3}\), and so on.

The partial sums converge. The sequence \(\frac{1}{1}, \frac{1}{2}, \frac{1}{3}, \dots\) decreases to \(0\), so the alternating series test shows that \(\lim_{n \to \infty} \sum_{k=1}^{n} x_k\) exists in \(\mathbb{R}\); in fact

\begin{equation*} \lim_{n \to \infty} \sum_{k=1}^{n} \frac{(-1)^k}{k} = -\ln 2 . \end{equation*}

(For an explicit argument that the limit exists, write \(s_n = \sum_{k=1}^{n} x_k\). The odd partial sums increase, since \(s_{2n+1} - s_{2n-1} = \frac{1}{2n} - \frac{1}{2n+1} > 0\); the even partial sums decrease, since \(s_{2n+2} - s_{2n} = -\frac{1}{2n+1} + \frac{1}{2n+2} < 0\); and \(s_{2n} - s_{2n+1} = \frac{1}{2n+1} > 0\), so every even partial sum exceeds every odd one. Thus \(s_1 \le s_3 \le \cdots\) is increasing and bounded above by \(s_2\), and \(s_2 \ge s_4 \ge \cdots\) is decreasing and bounded below by \(s_1\); both subsequences converge in \(\mathbb{R}\), and \(s_{2n} - s_{2n+1} = \frac{1}{2n+1} \to 0\) forces the two limits to agree. Hence \(\lim_{n \to \infty} s_n\) exists in \(\mathbb{R}\).)

The integral is not defined. Recall from 3.18 that \(\int x \, d\mu\) is defined only when at least one of \(\int x^+ \, d\mu\) and \(\int x^- \, d\mu\) is finite. Here

\begin{equation*} x^+(k) = \begin{cases} \dfrac{1}{k} & \text{if } k \text{ is even,} \\[4pt] 0 & \text{if } k \text{ is odd,} \end{cases} \qquad x^-(k) = \begin{cases} 0 & \text{if } k \text{ is even,} \\[4pt] \dfrac{1}{k} & \text{if } k \text{ is odd.} \end{cases} \end{equation*}

Both \(x^+\) and \(x^-\) are nonnegative functions on \(\mathbb{Z}^+\), and every subset of \(\mathbb{Z}^+\) belongs to the \(\sigma\)-algebra on which counting measure is defined, so both are measurable. By Example 3.6, integration with respect to counting measure on \(\mathbb{Z}^+\) is summation, so

\begin{equation*} \int x^+ \, d\mu = \sum_{k=1}^{\infty} x^+(k) = \sum_{m=1}^{\infty} \frac{1}{2m} = \frac{1}{2} \sum_{m=1}^{\infty} \frac{1}{m} = \infty \end{equation*}

and

\begin{equation*} \int x^- \, d\mu = \sum_{k=1}^{\infty} x^-(k) = \sum_{m=1}^{\infty} \frac{1}{2m - 1} \ge \sum_{m=1}^{\infty} \frac{1}{2m} = \infty , \end{equation*}

both because the harmonic series diverges. Since \(\int x^+ \, d\mu = \int x^- \, d\mu = \infty\), the expression \(\int x^+ \, d\mu - \int x^- \, d\mu\) is of the form \(\infty - \infty\), and hence \(\int x \, d\mu\) is not defined.

The point of the example is that \(\lim_{n \to \infty} \sum_{k=1}^{n} x_k\) takes the terms in a particular order and relies on cancellation, whereas integration with respect to counting measure is an order-free notion: it is defined only for series that converge absolutely (equivalently, only when \(\int |x| \, d\mu = \sum_{k=1}^{\infty} |x_k| < \infty\)) or for series in which only one of the two sign parts is infinite. Here \(\sum_{k=1}^{\infty} \frac{1}{k} = \infty\), so the convergence is merely conditional.

Problem (3A.19)

Show that if \((X, \mathcal{S}, \mu)\) is a measure space and \(f : X \to [0, \infty)\) is \(\mathcal{S}\)-measurable, then

\begin{equation*} \mu(X) \inf_X f \le \int f \, d\mu \le \mu(X) \sup_X f . \end{equation*}

Solution

Throughout we use the book’s convention, stated alongside 3.2, that \(0 \cdot \infty\) and \(\infty \cdot 0\) are interpreted to be \(0\). The quantities \(\inf_X f\) and \(\sup_X f\) presuppose \(X \ne \varnothing\), so assume that; then \(\inf_X f \in [0, \infty)\) and \(\sup_X f \in [0, \infty]\). (If \(X = \varnothing\) the assertion is vacuous: \(\mu(X) = 0\) and \(\int f \, d\mu = 0\).)

The left inequality. Let \(P\) be the \(\mathcal{S}\)-partition of \(X\) consisting of the single set \(A_1 = X\). (This is a legitimate \(\mathcal{S}\)-partition: it is a finite collection of pairwise disjoint sets in \(\mathcal{S}\) whose union is \(X\).) By definition 3.2 its lower Lebesgue sum is

\begin{equation*} L(f, P) = \mu(A_1) \inf_{A_1} f = \mu(X) \inf_X f . \end{equation*}

Since \(\int f \, d\mu\) is by definition 3.3 the supremum of \(L(f, Q)\) over all \(\mathcal{S}\)-partitions \(Q\) of \(X\), and \(P\) is one such partition,

\begin{equation*} \mu(X) \inf_X f = L(f, P) \le \int f \, d\mu . \end{equation*}

The right inequality. Let \(c = \sup_X f \in [0, \infty]\). Then \(f(x) \le c\) for every \(x \in X\), that is,

\begin{equation*} f(x) \le c \, \chi_X(x) \qquad \text{for every } x \in X . \end{equation*}

Both \(f\) and \(c \chi_X\) are \(\mathcal{S}\)-measurable functions from \(X\) to \([0, \infty]\), so 3.8 (integration is order preserving) gives

\begin{equation*} \int f \, d\mu \le \int c \, \chi_X \, d\mu . \end{equation*}

Applying 3.15 with \(n = 1\), \(E_1 = X \in \mathcal{S}\), and \(c_1 = c \in [0, \infty]\) yields

\begin{equation*} \int c \, \chi_X \, d\mu = c \, \mu(X) = \mu(X) \sup_X f . \end{equation*}

Combining the last two displays gives \(\int f \, d\mu \le \mu(X) \sup_X f\), which together with the left inequality completes the proof.

Two remarks on the degenerate cases. If \(\mu(X) = \infty\) and \(\inf_X f = 0\), the left inequality reads \(0 \le \int f \, d\mu\), which is true (and by Exercise 1 in this section, \(\inf_E f = 0\) is forced on every \(E\) of infinite measure whenever \(\int f \, d\mu < \infty\)). If \(\mu(X) = 0\), then for every \(\mathcal{S}\)-partition \(A_1, \dots, A_m\) of \(X\) we have \(\mu(A_j) \le \mu(X) = 0\) by 2.57, so each term \(\mu(A_j) \inf_{A_j} f\) is \(0 \cdot t\) with \(t \in [0, \infty]\) and hence equals \(0\) by the convention above; thus every lower Lebesgue sum vanishes, \(\int f \, d\mu = 0\), and both inequalities hold.

Problem (3A.20)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(f_1, f_2, \dots\) is a monotone (meaning either increasing or decreasing) sequence of \(\mathcal{S}\)-measurable functions. Define \(f : X \to [-\infty, \infty]\) by

\begin{equation*} f(x) = \lim_{k \to \infty} f_k(x). \end{equation*}

Prove that if \(\int |f_1| \, d\mu < \infty\), then

\begin{equation*} \lim_{k \to \infty} \int f_k \, d\mu = \int f \, d\mu . \end{equation*}

Solution

The functions \(f_k\) take values in \([-\infty, \infty]\); because the sequence is monotone, \(\lim_{k \to \infty} f_k(x)\) exists in \([-\infty, \infty]\) for each \(x\), so \(f\) is well defined, and \(f\) is \(\mathcal{S}\)-measurable by 2.53 (it is the pointwise supremum of the \(f_k\) in the increasing case and the pointwise infimum in the decreasing case).

We first record two lemmas, then treat the increasing case, then deduce the decreasing case.

Lemma A. If \(N \in \mathcal{S}\) with \(\mu(N) = 0\) and \(h : X \to [0, \infty]\) is \(\mathcal{S}\)-measurable, then \(h \chi_N\) is \(\mathcal{S}\)-measurable and \(\int h \chi_N \, d\mu = 0\).

Proof of Lemma A. For \(a \ge 0\) we have \((h \chi_N)^{-1}\bigl((a, \infty]\bigr) = h^{-1}\bigl((a,\infty]\bigr) \cap N \in \mathcal{S}\), and for \(a < 0\) we have \((h \chi_N)^{-1}\bigl((a, \infty]\bigr) = X \in \mathcal{S}\); hence \(h \chi_N\) is \(\mathcal{S}\)-measurable by 2.52. Now let \(P\) be any \(\mathcal{S}\)-partition \(A_1, \dots, A_m\) of \(X\) and fix \(j\). If \(A_j \not\subseteq N\), choose \(y \in A_j \setminus N\); then \((h\chi_N)(y) = 0\), so \(\inf_{A_j} h \chi_N = 0\) and therefore \(\mu(A_j) \inf_{A_j} h \chi_N = 0\) (using the convention \(\infty \cdot 0 = 0\)). If \(A_j \subseteq N\), then \(\mu(A_j) \le \mu(N) = 0\) by 2.57, so again \(\mu(A_j) \inf_{A_j} h \chi_N = 0\) (using \(0 \cdot \infty = 0\)). Hence \(L(h\chi_N, P) = 0\) for every \(\mathcal{S}\)-partition \(P\), and definition 3.3 gives \(\int h \chi_N \, d\mu = 0\).

Lemma B. Suppose \(N \in \mathcal{S}\) with \(\mu(N) = 0\), put \(M = X \setminus N\), and suppose \(g : X \to [-\infty,\infty]\) is \(\mathcal{S}\)-measurable and \(\int g \, d\mu\) is defined. Then \(g \chi_M\) is \(\mathcal{S}\)-measurable, \(\int g \chi_M \, d\mu\) is defined, and \(\int g \chi_M \, d\mu = \int g \, d\mu\).

Proof of Lemma B. Measurability: for \(a \ge 0\), \((g\chi_M)^{-1}\bigl((a,\infty]\bigr) = g^{-1}\bigl((a,\infty]\bigr) \cap M \in \mathcal{S}\), and for \(a < 0\), \((g \chi_M)^{-1}\bigl((a,\infty]\bigr) = \bigl(g^{-1}((a,\infty]) \cap M\bigr) \cup N \in \mathcal{S}\); apply 2.52. Next, \((g\chi_M)^+ = g^+ \chi_M\) and \((g \chi_M)^- = g^- \chi_M\) pointwise. Since \(g^{\pm} = g^{\pm} \chi_M + g^{\pm} \chi_N\) pointwise, additivity 3.16 and Lemma A give

\begin{equation*} \int g^{\pm} \, d\mu = \int g^{\pm} \chi_M \, d\mu + \int g^{\pm} \chi_N \, d\mu = \int (g \chi_M)^{\pm} \, d\mu . \end{equation*}

Thus \(\int (g\chi_M)^+ d\mu\) and \(\int (g \chi_M)^- d\mu\) equal \(\int g^+ d\mu\) and \(\int g^- d\mu\) respectively; in particular at least one is finite, so \(\int g\chi_M \, d\mu\) is defined and equals \(\int g \, d\mu\).

The increasing case. Assume \(f_1 \le f_2 \le \cdots\).

Step 0: all the integrals are defined. For each \(k\) we have \(f_k \ge f_1\), hence \(f_k^- \le f_1^- \le |f_1|\) pointwise, so 3.8 gives \(\int f_k^- \, d\mu \le \int |f_1| \, d\mu < \infty\). Thus \(\int f_k \, d\mu\) is defined (and lies in \((-\infty, \infty]\)). The same argument with \(f\) in place of \(f_k\), using \(f \ge f_1\), shows \(\int f \, d\mu\) is defined.

Step 1: the set where \(f_1\) is infinite is null. Let \(N = f_1^{-1}\bigl(\{-\infty, \infty\}\bigr)\), which lies in \(\mathcal{S}\) because \(f_1\) is \(\mathcal{S}\)-measurable and \(\{-\infty, \infty\}\) is a Borel subset of \([-\infty, \infty]\) (its intersection with \(\mathbb{R}\) is \(\varnothing\); see 2.50). For every \(n \in \mathbb{Z}^+\) we have \(|f_1| \ge n \chi_N\) pointwise, so 3.8 together with 3.15 (with \(n = 1\) summand, \(E_1 = N\), \(c_1 = n\)) gives \(\int |f_1| \, d\mu \ge n \mu(N)\). Since \(\int |f_1| \, d\mu < \infty\) and \(n\) is arbitrary, \(\mu(N) = 0\).

Step 2: reduction to the case where \(f_1\) is real valued. Let \(M = X \setminus N\) and set \(\widetilde{f_k} = f_k \chi_M\) and \(\widetilde{f} = f \chi_M\). By Lemma B these are \(\mathcal{S}\)-measurable and

\begin{equation*} \int \widetilde{f_k} \, d\mu = \int f_k \, d\mu \quad (k \in \mathbb{Z}^+), \qquad \int \widetilde{f} \, d\mu = \int f \, d\mu . \end{equation*}

On \(N\) every \(\widetilde{f_k}\) and \(\widetilde{f}\) equals \(0\); on \(M\) they equal \(f_k\) and \(f\). Hence \(\widetilde{f_1} \le \widetilde{f_2} \le \cdots\) is still increasing, \(\lim_{k \to \infty} \widetilde{f_k}(x) = \widetilde{f}(x)\) for every \(x \in X\), and \(|\widetilde{f_1}| \le |f_1|\) gives \(\int |\widetilde{f_1}| \, d\mu < \infty\) by 3.8. Moreover \(\widetilde{f_1}\) is real valued. So it suffices to prove the result under the additional hypothesis that \(f_1(x) \in \mathbb{R}\) for every \(x \in X\), which we now assume.

Step 3: the auxiliary nonnegative sequence. Because \(f_k \ge f_1 > -\infty\), each \(f_k\) takes values in \((-\infty, \infty]\), and likewise \(f\) does. Since \(f_1\) is real valued, the differences

\begin{equation*} g_k := f_k - f_1 \quad \text{and} \quad g := f - f_1 \end{equation*}

are well defined functions from \(X\) to \([0, \infty]\). They are \(\mathcal{S}\)-measurable: for \(a \in \mathbb{R}\),

\begin{equation*} g_k^{-1}\bigl((a, \infty]\bigr) = \{x : f_k(x) > a + f_1(x)\} = \bigcup_{r \in \mathbb{Q}} \Bigl( f_k^{-1}\bigl((r, \infty]\bigr) \cap f_1^{-1}\bigl([-\infty, r - a)\bigr) \Bigr) , \end{equation*}

because \(f_k(x) > a + f_1(x)\) (with \(a + f_1(x) \in \mathbb{R}\)) holds if and only if there is a rational \(r\) with \(a + f_1(x) < r < f_k(x)\), i.e. \(f_k(x) > r\) and \(f_1(x) < r - a\). The right side is a countable union of sets in \(\mathcal{S}\), so \(g_k\) is \(\mathcal{S}\)-measurable by 2.52; the same computation applies to \(g\). Clearly \(0 \le g_1 \le g_2 \le \cdots\) and \(\lim_{k \to \infty} g_k(x) = g(x)\) for every \(x \in X\). Hence the Monotone Convergence Theorem 3.11 gives

\begin{equation*} \lim_{k \to \infty} \int g_k \, d\mu = \int g \, d\mu . \end{equation*}

Step 4: splitting off \(f_1\). We claim that if \(h : X \to (-\infty, \infty]\) is \(\mathcal{S}\)-measurable with \(h \ge f_1\), then

\begin{equation*} \int h \, d\mu = \int (h - f_1) \, d\mu + \int f_1 \, d\mu . \end{equation*}

First, \(h - f_1\) is a well-defined \(\mathcal{S}\)-measurable function from \(X\) to \([0, \infty]\), by the rational-splitting computation of Step 3 with \(h\) in place of \(f_k\) (that computation used only that \(f_1\) is real valued, that \(h\) is \(\mathcal{S}\)-measurable, and that \(h \ge f_1\)). Next, note the pointwise identity between functions with values in \([0, \infty]\)

\begin{equation*} h^+ + f_1^- = (h - f_1) + f_1^+ + h^- , \end{equation*}

which is verified by substituting \(h - f_1 = h^+ - h^- - f_1^+ + f_1^-\) (legitimate at every \(x\), since \(f_1(x)\) is real and at most one of \(h^+(x), h^-(x)\) is nonzero) and cancelling; when \(h(x) = \infty\) both sides equal \(\infty\). Integrating both sides and using additivity 3.16 twice,

\begin{equation*} \int h^+ \, d\mu + \int f_1^- \, d\mu = \int (h - f_1) \, d\mu + \int f_1^+ \, d\mu + \int h^- \, d\mu . \end{equation*}

Now \(\int f_1^+ \, d\mu\) and \(\int f_1^- \, d\mu\) are finite because \(f_1^{\pm} \le |f_1|\) and \(\int |f_1| \, d\mu < \infty\), and \(\int h^- \, d\mu \le \int f_1^- \, d\mu < \infty\) because \(h \ge f_1\) forces \(h^- \le f_1^-\). Subtracting the finite numbers \(\int f_1^- \, d\mu\) and \(\int h^- \, d\mu\) from both sides gives

\begin{equation*} \int h^+ \, d\mu - \int h^- \, d\mu = \int (h - f_1) \, d\mu + \int f_1^+ \, d\mu - \int f_1^- \, d\mu , \end{equation*}

which by 3.18 is exactly the claim (both sides possibly equal to \(+\infty\)).

Step 5: conclusion in the increasing case. Applying Step 4 with \(h = f_k\) and with \(h = f\),

\begin{equation*} \int f_k \, d\mu = \int g_k \, d\mu + \int f_1 \, d\mu \quad \text{and} \quad \int f \, d\mu = \int g \, d\mu + \int f_1 \, d\mu . \end{equation*}

Since \(\int f_1 \, d\mu\) is a real number, adding it commutes with taking limits in \((-\infty, \infty]\), so Step 3 gives

\begin{equation*} \lim_{k \to \infty} \int f_k \, d\mu = \lim_{k \to \infty} \int g_k \, d\mu + \int f_1 \, d\mu = \int g \, d\mu + \int f_1 \, d\mu = \int f \, d\mu . \end{equation*}

The decreasing case. Now assume \(f_1 \ge f_2 \ge \cdots\). Put \(h_k = -f_k\). Each \(h_k\) is \(\mathcal{S}\)-measurable, since \(h_k^{-1}\bigl((a, \infty]\bigr) = f_k^{-1}\bigl([-\infty, -a)\bigr) \in \mathcal{S}\) for every \(a \in \mathbb{R}\) (use 2.52). Then \(h_1 \le h_2 \le \cdots\) is increasing, \(\lim_{k \to \infty} h_k(x) = -f(x)\) for every \(x\), and \(\int |h_1| \, d\mu = \int |f_1| \, d\mu < \infty\). Also \(\int f_k \, d\mu\) is defined for each \(k\): from \(f_k \le f_1\) we get \(f_k^+ \le f_1^+ \le |f_1|\), so \(\int f_k^+ \, d\mu < \infty\) by 3.8; similarly \(\int f \, d\mu\) is defined. By the increasing case already proved,

\begin{equation*} \lim_{k \to \infty} \int (-f_k) \, d\mu = \int (-f) \, d\mu . \end{equation*}

By homogeneity 3.20 with \(c = -1\) (applicable because \(\int f_k \, d\mu\) and \(\int f \, d\mu\) are defined), \(\int (-f_k) \, d\mu = -\int f_k \, d\mu\) and \(\int(-f) \, d\mu = -\int f \, d\mu\). Substituting and multiplying by \(-1\) yields

\begin{equation*} \lim_{k \to \infty} \int f_k \, d\mu = \int f \, d\mu , \end{equation*}

completing the proof.

Remark. Exercise 14 shows that the conclusion can fail for a decreasing sequence without some integrability hypothesis (for instance \(f_k = \chi_{[k, \infty)}\) with Lebesgue measure on \(\mathbb{R}\), where \(\int f_k \, d\lambda = \infty\) for every \(k\) but \(f = 0\)). The hypothesis \(\int |f_1| \, d\mu < \infty\) is exactly what rules this out: it pins the whole decreasing sequence below an integrable function. This is why the note after Exercise 14 points to the present exercise.

Problem (3A.21)

Henri Lebesgue wrote the following about his method of integration:

I have to pay a certain sum, which I have collected in my pocket. I take the bills and coins out of my pocket and give them to the creditor in the order I find them until I have reached the total sum. This is the Riemann integral. But I can proceed differently. After I have taken all the money out of my pocket I order the bills and coins according to identical values and then I pay the several heaps one after the other to the creditor. This is my integral.

Use 3.15 to explain what Lebesgue meant and to explain why integration of a function with respect to a measure can be thought of as partitioning the range of the function, in contrast to Riemann integration, which depends on partitioning the domain of the function.

[The quote above is taken from page 796 of The Princeton Companion to Mathematics, edited by Timothy Gowers.]

Solution

Recall the statement of 3.15: if \((X, \mathcal{S}, \mu)\) is a measure space, \(E_1, \dots, E_n \in \mathcal{S}\), and \(c_1, \dots, c_n \in [0, \infty]\), then

\begin{equation*} \int \sum_{k=1}^{n} c_k \chi_{E_k} \, d\mu = \sum_{k=1}^{n} c_k \, \mu(E_k). \end{equation*}

Everything in Lebesgue’s parable is contained in this formula.

The pocket, literally. Model the money as follows. Let \(X\) be the (finite) set of bills and coins in the pocket, let \(\mathcal{S}\) be the \(\sigma\)-algebra of all subsets of \(X\), and let \(\mu\) be counting measure on \(\mathcal{S}\), so that \(\mu(E)\) is the number of items in \(E\). Let \(v : X \to [0, \infty)\) assign to each item its face value. The total to be paid is \(\int v \, d\mu\): indeed, writing \(v = \sum_{x \in X} v(x) \chi_{\{x\}}\) (a finite sum, since \(X\) is finite) and applying 3.15 with \(\mu(\{x\}) = 1\) gives \(\int v \, d\mu = \sum_{x \in X} v(x)\), literally the sum of the values of all the items in the pocket. (This is the finite case of Example 3.6, where integration against counting measure is summation.)

Lebesgue’s method is 3.15 read from left to right. Let \(c_1, \dots, c_n\) be the distinct denominations present, and let

\begin{equation*} E_k = v^{-1}(\{c_k\}) \end{equation*}

be the heap of all items of value \(c_k\). These heaps are exactly what Lebesgue means by ordering the money “according to identical values”: each \(E_k\) is a piece of the domain cut out by a single value in the range. Since \(v = \sum_{k=1}^{n} c_k \chi_{E_k}\), formula 3.15 gives

\begin{equation*} \int v \, d\mu = \sum_{k=1}^{n} c_k \, \mu(E_k) = \sum_{k=1}^{n} (\text{denomination } c_k) \times (\text{number of items in the heap}) . \end{equation*}

That is precisely “paying the several heaps one after the other”: you hand over the pile of tens, then the pile of fives, and so on, and the sum is a sum over the \(n\) values in the range, not over the many items in the domain. Riemann’s method, by contrast, is “in the order I find them”: you fix an enumeration of the items in the pocket, i.e. an ordering of the domain, and accumulate the values one item at a time, running along the domain.

Why 3.15 is a range partition. Note what 3.15 asks of the sets \(E_1, \dots, E_n\): only that they belong to \(\mathcal{S}\), so that the numbers \(\mu(E_k)\) make sense. They need not be intervals, and 3.15 does not even require them to be disjoint. The numbers \(c_1, \dots, c_n\), on the other hand, are values of the function. So the data of the sum \(\sum_k c_k \mu(E_k)\) is: a list of values, and the measure of the set on which each value is taken.

This is exactly how the integral of an arbitrary nonnegative measurable function is computed. Let \(f : X \to [0, \infty)\) be \(\mathcal{S}\)-measurable. For each \(j \in \mathbb{Z}^+\) chop the part \([0, j)\) of the range into the \(j 2^{j}\) half-open intervals \([k 2^{-j}, (k+1) 2^{-j})\) with \(0 \le k < j 2^{j}\), lump the rest of the range together, and let

\begin{equation*} E_{j,k} = f^{-1}\bigl([k 2^{-j}, (k+1) 2^{-j})\bigr), \qquad E_{j,\infty} = f^{-1}\bigl([j, \infty)\bigr) . \end{equation*}

These sets lie in \(\mathcal{S}\) precisely because \(f\) is \(\mathcal{S}\)-measurable, and for fixed \(j\) they are disjoint with union \(X\). Put

\begin{equation*} s_j = \sum_{k=0}^{j 2^{j} - 1} k 2^{-j} \, \chi_{E_{j,k}} + j \, \chi_{E_{j, \infty}} , \end{equation*}

a finite linear combination of characteristic functions, so that 3.15 applies verbatim and gives

\begin{equation*} \int s_j \, d\mu = \sum_{k=0}^{j 2^{j} - 1} k 2^{-j} \, \mu(E_{j,k}) + j \, \mu(E_{j,\infty}) . \end{equation*}

Pointwise \(s_j = \min\{ 2^{-j} \lfloor 2^{j} f \rfloor, \, j \}\); since \(2^{-j} \lfloor 2^{j} t \rfloor\) increases with \(j\) for each fixed \(t \ge 0\), we get \(0 \le s_1 \le s_2 \le \cdots\), and \(s_j(x) \to f(x)\) for every \(x\) (once \(j > f(x)\) we have \(0 \le f(x) - s_j(x) \le 2^{-j}\)). This is the construction behind 2.89 in the case of a nonnegative function, and the Monotone Convergence Theorem 3.11 now gives \(\int f \, d\mu\) as the limit of the sums displayed above. The partition being refined at each stage is a partition of the range; the domain pieces \(E_{j,k}\) are whatever the function’s level sets happen to be, and the only thing demanded of them is that they be measurable.

Contrast with Riemann integration. A Riemann sum for \(f\) on \([a, b]\) starts from a partition \(a = t_0 < t_1 < \cdots < t_n = b\) of the domain and forms \(\sum_j f(s_j)(t_j - t_{j-1})\) with \(s_j \in [t_{j-1}, t_j]\); the lower Riemann sums \(L(f, P, [a,b])\) of 1.3 replace \(f(s_j)\) by \(\inf_{[t_{j-1}, t_j]} f\). The accuracy of such a sum is governed by the oscillation of \(f\) on each subinterval, so the method works only when the function is nearly constant on small intervals of the domain. Nothing about the range is used to choose the partition.

The two methods therefore fail and succeed in different places. For \(f = \chi_{\mathbb{Q} \cap [0,1]}\) every subinterval of \([0,1]\) contains both rationals and irrationals, so \(\inf f = 0\) and \(\sup f = 1\) on every piece of every domain partition; hence \(L(f, [0,1]) = 0\) while \(U(f, [0,1]) = 1\), and \(f\) is not Riemann integrable. But the range of \(f\) has only the two points \(0\) and \(1\), so Lebesgue’s sorting produces just two heaps, \(\mathbb{Q} \cap [0,1]\) and \([0,1] \setminus \mathbb{Q}\), and 3.15 gives at once

\begin{equation*} \int \chi_{\mathbb{Q} \cap [0,1]} \, d\lambda = 1 \cdot \lambda\bigl(\mathbb{Q} \cap [0,1]\bigr) + 0 \cdot \lambda\bigl([0,1] \setminus \mathbb{Q}\bigr) = 1 \cdot 0 + 0 \cdot 1 = 0 , \end{equation*}

in agreement with 3.5. A wildly oscillating function still has a small range, and grouping the domain by value tames it.

The price of Lebesgue’s method is visible in the same formula: to pay the heaps one must be able to count each heap, that is, one must know \(\mu(E_k)\) for level sets \(E_k\) that may be very complicated sets. This is exactly why a theory of measure (Chapter 2) has to be built before the integral can be defined, and it is why the definition 3.3 of \(\int f \, d\mu\) via lower Lebesgue sums allows arbitrary sets of \(\mathcal{S}\) in its partitions rather than only intervals, as the lower Riemann sums do. Once arbitrary measurable sets are allowed, a domain partition can be chosen to follow the level sets of \(f\), and the domain-partition definition 3.3 becomes equivalent to the range-partition picture described above.

Exercises 3B

Problem (3B.1)

Give an example of a sequence \(f_1, f_2, \ldots\) of functions from \(\mathbb{Z}^+\) to \([0, \infty)\) such that

\begin{equation*} \lim_{k \to \infty} f_k(m) = 0 \end{equation*}

for every \(m \in \mathbb{Z}^+\) but \(\lim_{k \to \infty} \int f_k \, d\mu = 1\), where \(\mu\) is counting measure on \(\mathbb{Z}^+\).

Solution

Take

\begin{equation*} f_k = \chi_{\{k\}}, \end{equation*}

that is, \(f_k : \mathbb{Z}^+ \to [0, \infty)\) is defined by \(f_k(m) = 1\) if \(m = k\) and \(f_k(m) = 0\) if \(m \neq k\). Here the \(\sigma\)-algebra is the collection of all subsets of \(\mathbb{Z}^+\), so every function on \(\mathbb{Z}^+\) is measurable.

Pointwise convergence to \(0\): fix \(m \in \mathbb{Z}^+\). If \(k > m\) then \(k \neq m\) and hence \(f_k(m) = 0\). Thus \(f_k(m) = 0\) for all \(k > m\), and in particular

\begin{equation*} \lim_{k \to \infty} f_k(m) = 0 . \end{equation*}

The integrals: because \(f_k\) is the characteristic function of the set \(\{k\}\), the formula for the integral of a characteristic function (3.4) gives

\begin{equation*} \int f_k \, d\mu = \mu(\{k\}) = 1 \end{equation*}

for every \(k \in \mathbb{Z}^+\), since counting measure assigns to each singleton the value \(1\). Hence

\begin{equation*} \lim_{k \to \infty} \int f_k \, d\mu = 1 \neq 0 = \int \Bigl( \lim_{k \to \infty} f_k \Bigr) d\mu . \end{equation*}

Thus the sequence \(f_1, f_2, \ldots\) has the required properties. Note that this example is consistent with the Dominated Convergence Theorem (3.31): the smallest function dominating all the \(f_k\) is \(g \equiv 1\), and \(\int g \, d\mu = \mu(\mathbb{Z}^+) = \infty\), so no dominating function with finite integral exists. The example is also consistent with the Monotone Convergence Theorem (3.11), since the sequence \(f_1, f_2, \ldots\) is not increasing. The mass \(1\) carried by each \(f_k\) simply escapes to infinity.

Problem (3B.2)

Give an example of a sequence \(f_1, f_2, \ldots\) of continuous functions from \(\mathbb{R}\) to \([0, 1]\) such that

\begin{equation*} \lim_{k \to \infty} f_k(x) = 0 \end{equation*}

for every \(x \in \mathbb{R}\) but \(\lim_{k \to \infty} \int f_k \, d\lambda = \infty\), where \(\lambda\) is Lebesgue measure on \(\mathbb{R}\).

Solution

For \(k \in \mathbb{Z}^+\) define \(f_k : \mathbb{R} \to [0,1]\) by

\begin{equation*} f_k(x) = \max\Bigl\{ 0, \; 1 - \frac{|x - k^2|}{k} \Bigr\} . \end{equation*}

Thus \(f_k\) is the tent function that equals \(1\) at \(x = k^2\), decreases linearly to \(0\) at the two points \(k^2 - k\) and \(k^2 + k\), and equals \(0\) outside the interval \([k^2 - k, \, k^2 + k]\).

Each \(f_k\) is continuous, since \(x \mapsto 1 - |x - k^2| / k\) is continuous and the maximum of two continuous functions is continuous. Clearly \(0 \le f_k(x) \le 1\) for all \(x\), so \(f_k\) maps \(\mathbb{R}\) into \([0,1]\). Being continuous, each \(f_k\) is Borel measurable (2.41).

Pointwise convergence to \(0\): fix \(x \in \mathbb{R}\). Because \(k^2 - k \to \infty\) as \(k \to \infty\), there exists \(K \in \mathbb{Z}^+\) such that \(k^2 - k > |x|\) for all \(k \ge K\). For such \(k\) we have \(x < k^2 - k\), hence \(x \notin [k^2 - k, \, k^2 + k]\) and therefore \(f_k(x) = 0\). Thus

\begin{equation*} \lim_{k \to \infty} f_k(x) = 0 \end{equation*}

for every \(x \in \mathbb{R}\).

The integrals: \(f_k\) vanishes outside \([k^2 - k, \, k^2 + k]\) and is continuous, hence Riemann integrable on that interval; by 3.34 its Lebesgue integral over the interval equals its Riemann integral there, which is the area of a triangle with base of length \(2k\) and height \(1\). Thus

\begin{equation*} \int f_k \, d\lambda = \int_{[k^2 - k, \, k^2 + k]} f_k \, d\lambda = \frac{1}{2} \cdot 2k \cdot 1 = k . \end{equation*}

(The integral over the complement of \([k^2-k,\,k^2+k]\) is \(0\) because \(f_k\) is \(0\) there.) Hence

\begin{equation*} \lim_{k \to \infty} \int f_k \, d\lambda = \lim_{k \to \infty} k = \infty , \end{equation*}

while \(\int (\lim_{k \to \infty} f_k) \, d\lambda = \int 0 \, d\lambda = 0\).

As in the previous exercise, the point is that the bumps run off to infinity while growing wider. Note that the Bounded Convergence Theorem (3.26) does not apply because \(\lambda(\mathbb{R}) = \infty\), and the Dominated Convergence Theorem (3.31) does not apply because any \(g\) with \(g \ge f_k\) for all \(k\) satisfies \(\int g \, d\lambda \ge \int f_k \, d\lambda = k\) for every \(k\), so \(\int g \, d\lambda = \infty\).

Problem (3B.3)

Suppose \(\lambda\) is Lebesgue measure on \(\mathbb{R}\) and \(f : \mathbb{R} \to \mathbb{R}\) is a Borel measurable function such that \(\int |f| \, d\lambda < \infty\). Define \(g : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} g(x) = \int_{(-\infty, \, x)} f \, d\lambda . \end{equation*}

Prove that \(g\) is uniformly continuous on \(\mathbb{R}\).

Solution

First we check that \(g\) is well defined and real valued. Fix \(x \in \mathbb{R}\). By 3.24 the expression \(\int_{(-\infty, x)} f \, d\lambda\) means \(\int \chi_{(-\infty, x)} f \, d\lambda\). Since \(|\chi_{(-\infty, x)} f| \le |f|\) pointwise, 3.8 gives

\begin{equation*} \int \bigl| \chi_{(-\infty, x)} f \bigr| \, d\lambda \le \int |f| \, d\lambda < \infty , \end{equation*}

so both \((\chi_{(-\infty,x)} f)^+\) and \((\chi_{(-\infty,x)} f)^-\) have finite integrals; hence the integral defining \(g(x)\) is defined (3.18) and is a real number. Also \(\chi_{(-\infty,x)} f\) is Borel measurable, being the product of a Borel measurable function and the characteristic function of a Borel set.

Now let \(\varepsilon > 0\). Apply 3.28 (integrals on small sets are small) to the nonnegative Borel measurable function \(|f|\), whose integral is finite: there exists \(\delta > 0\) such that

\begin{equation*} \int_B |f| \, d\lambda < \varepsilon \quad \text{whenever } B \text{ is a Borel set with } \lambda(B) < \delta . \end{equation*}

Suppose now that \(x, y \in \mathbb{R}\) with \(|x - y| < \delta\); without loss of generality \(x \le y\). Then

\begin{equation*} \chi_{(-\infty, \, y)} f = \chi_{(-\infty, \, x)} f + \chi_{[x, \, y)} f , \end{equation*}

because the intervals \((-\infty, x)\) and \([x,y)\) are disjoint with union \((-\infty, y)\). Both functions on the right have absolute value bounded by \(|f|\) and hence have finite integrals of their absolute values, so additivity of integration (3.21) applies and gives

\begin{equation*} g(y) - g(x) = \int \chi_{[x, \, y)} f \, d\lambda = \int_{[x, \, y)} f \, d\lambda . \end{equation*}

Therefore, using 3.23 and then 3.8,

\begin{equation*} |g(y) - g(x)| = \left| \int_{[x, \, y)} f \, d\lambda \right| \le \int_{[x, \, y)} |f| \, d\lambda . \end{equation*}

Since \(\lambda([x,y)) = y - x = |x - y| < \delta\), the choice of \(\delta\) gives

\begin{equation*} \int_{[x, \, y)} |f| \, d\lambda < \varepsilon . \end{equation*}

Hence \(|g(y) - g(x)| < \varepsilon\).

We have shown: for every \(\varepsilon > 0\) there exists \(\delta > 0\) such that \(|g(x) - g(y)| < \varepsilon\) for all \(x, y \in \mathbb{R}\) with \(|x - y| < \delta\). Because \(\delta\) depends only on \(\varepsilon\) and not on the points \(x, y\), this is exactly the statement that \(g\) is uniformly continuous on \(\mathbb{R}\).

Problem (3B.4)

(a) Suppose \((X, \mathcal{S}, \mu)\) is a measure space with \(\mu(X) < \infty\). Suppose that \(f : X \to [0, \infty)\) is a bounded \(\mathcal{S}\)-measurable function. Prove that

\begin{equation*} \int f \, d\mu = \inf\Bigl\{ \sum_{j=1}^{m} \mu(A_j) \sup_{A_j} f \; : \; A_1, \ldots, A_m \text{ is an } \mathcal{S}\text{-partition of } X \Bigr\} . \end{equation*}

(b) Show that the conclusion of (a) can fail if the hypothesis that \(f\) is bounded is replaced by the hypothesis that \(\int f \, d\mu < \infty\).

(c) Show that the conclusion of (a) can fail if the condition that \(\mu(X) < \infty\) is deleted.

[Part (a) of this exercise shows that if we had defined an upper Lebesgue sum, then it could be used to define \(\int f \, d\mu\) when \(f\) is bounded and \(\mu(X) < \infty\). However, parts (b) and (c) show that the hypotheses that \(f\) is bounded and that \(\mu(X) < \infty\) are needed if defining the integral via the equation above. The definition of the integral via the lower Lebesgue sum does not require these hypotheses, showing the advantage of using the lower Lebesgue sum.]

Solution

Throughout, for an \(\mathcal{S}\)-partition \(P\) of \(X\) given by \(A_1, \ldots, A_m\) we write, in parallel with the lower Lebesgue sum \(L(f, P) = \sum_{j=1}^m \mu(A_j) \inf_{A_j} f\) of 3.2, the upper Lebesgue sum

\begin{equation*} U(f, P) = \sum_{j=1}^{m} \mu(A_j) \sup_{A_j} f , \end{equation*}

and we set \(U(f) = \inf \{ U(f, P) : P \text{ is an } \mathcal{S}\text{-partition of } X \}\). The assertion of (a) is that \(\int f \, d\mu = U(f)\).

Part (a).

If \(X = \varnothing\) then \(\mu(X) = 0\) and both sides equal \(0\), so assume \(X \neq \varnothing\). The definition 3.1 of an \(\mathcal{S}\)-partition permits empty sets, but a term of \(L(f, P)\) or of \(U(f, P)\) coming from an empty \(A_j\) has the factor \(\mu(\varnothing) = 0\), hence is \(0\) whatever value is assigned to \(\inf_{\varnothing} f\) or \(\sup_{\varnothing} f\), by the convention (stated just before 3.2) that \(0 \cdot \infty = \infty \cdot 0 = 0\). Deleting the empty sets from a partition of the nonempty set \(X\) therefore leaves an \(\mathcal{S}\)-partition of \(X\) with the same lower and upper sums. So we may assume throughout that all sets appearing in a partition are nonempty; then \(\inf_{A_j} f\) and \(\sup_{A_j} f\) are genuine real numbers, because \(f\) is bounded. Let \(c \in [0, \infty)\) be such that \(0 \le f(x) \le c\) for all \(x \in X\). Since \(\mu(X) < \infty\), every \(L(f, P)\) and every \(U(f, P)\) is a real number in \([0, c\,\mu(X)]\).

Step 1: refining a partition increases \(L\) and decreases \(U\). Let \(P\) be the partition \(A_1, \ldots, A_m\) and let \(P^{\prime}\) be the partition \(B_1, \ldots, B_n\). Let \(Q\) be the \(\mathcal{S}\)-partition consisting of the nonempty sets among \(A_i \cap B_j\) for \(i \in \{1, \ldots, m\}\) and \(j \in \{1, \ldots, n\}\); these sets are disjoint, lie in \(\mathcal{S}\), and their union is \(X\). Fix \(i\). The sets \(A_i \cap B_1, \ldots, A_i \cap B_n\) are disjoint with union \(A_i\), so additivity of the measure \(\mu\) gives \(\mu(A_i) = \sum_{j=1}^n \mu(A_i \cap B_j)\). Moreover \(\inf_{A_i} f \le \inf_{A_i \cap B_j} f\) whenever \(A_i \cap B_j \neq \varnothing\), because the infimum is taken over a smaller set. Hence

\begin{equation*} \mu(A_i) \inf_{A_i} f = \sum_{j=1}^{n} \mu(A_i \cap B_j) \inf_{A_i} f \le \sum_{j=1}^{n} \mu(A_i \cap B_j) \inf_{A_i \cap B_j} f , \end{equation*}

where terms with \(A_i \cap B_j = \varnothing\) contribute \(0\) on both sides. Summing over \(i \in \{1, \ldots, m\}\) yields \(L(f, P) \le L(f, Q)\). The same computation with \(B_j\) in place of \(A_i\) and suprema in place of infima (now using \(\sup_{A_i \cap B_j} f \le \sup_{B_j} f\)) yields \(U(f, Q) \le U(f, P^{\prime})\).

Step 2: \(\int f \, d\mu \le U(f)\). Since \(\inf_{A} f \le \sup_{A} f\) for every nonempty \(A\), we have \(L(f, Q) \le U(f, Q)\) for every partition \(Q\). Combining this with Step 1, for any two \(\mathcal{S}\)-partitions \(P\) and \(P^{\prime}\) we get

\begin{equation*} L(f, P) \le L(f, Q) \le U(f, Q) \le U(f, P^{\prime}) , \end{equation*}

where \(Q\) is the common refinement constructed above. Fix \(P^{\prime}\) and take the supremum over all \(P\): by the definition 3.3 of the integral this gives \(\int f \, d\mu \le U(f, P^{\prime})\). Now take the infimum over all \(P^{\prime}\) to obtain \(\int f \, d\mu \le U(f)\).

Step 3: \(U(f) \le \int f \, d\mu\). If \(c = 0\) then \(f \equiv 0\) and both quantities are \(0\); so assume \(c > 0\). Let \(\varepsilon > 0\) and choose \(n \in \mathbb{Z}^+\) with \(\frac{c}{n} \mu(X) < \varepsilon\) (possible because \(\mu(X) < \infty\)). Define

\begin{equation*} A_j = f^{-1}\Bigl( \Bigl[ \tfrac{(j-1)c}{n}, \tfrac{jc}{n} \Bigr) \Bigr) \quad \text{for } j = 1, \ldots, n-1, \qquad A_n = f^{-1}\Bigl( \Bigl[ \tfrac{(n-1)c}{n}, c \Bigr] \Bigr) . \end{equation*}

Each \(A_j\) is in \(\mathcal{S}\) because \(f\) is \(\mathcal{S}\)-measurable and the sets used are Borel subsets of \(\mathbb{R}\); the sets \(A_1, \ldots, A_n\) are disjoint, and their union is \(X\) because \(f(X) \subseteq [0, c]\). So \(P\), the list of nonempty sets among \(A_1, \ldots, A_n\), is an \(\mathcal{S}\)-partition of \(X\). For each nonempty \(A_j\) the values of \(f\) on \(A_j\) lie in an interval of length \(c/n\), so

\begin{equation*} \sup_{A_j} f - \inf_{A_j} f \le \frac{c}{n} . \end{equation*}

Therefore

\begin{equation*} U(f, P) - L(f, P) = \sum_{j} \mu(A_j) \Bigl( \sup_{A_j} f - \inf_{A_j} f \Bigr) \le \frac{c}{n} \sum_{j} \mu(A_j) = \frac{c}{n} \mu(X) < \varepsilon , \end{equation*}

this subtraction being legitimate because all the numbers involved are finite. Hence

\begin{equation*} U(f) \le U(f, P) < L(f, P) + \varepsilon \le \int f \, d\mu + \varepsilon , \end{equation*}

the last inequality holding by 3.3. Since \(\varepsilon > 0\) was arbitrary, \(U(f) \le \int f \, d\mu\).

Steps 2 and 3 together give \(\int f \, d\mu = U(f)\), which is the desired equality.

Part (b).

Let \(X = \mathbb{Z}^+\), let \(\mathcal{S}\) be the \(\sigma\)-algebra of all subsets of \(\mathbb{Z}^+\), and define \(\mu\) on \(\mathcal{S}\) by

\begin{equation*} \mu(A) = \sum_{n \in A} 2^{-n} . \end{equation*}

Then \(\mu(\varnothing) = 0\), and \(\mu\) is countably additive because a series of nonnegative terms may be summed over a countable disjoint decomposition of the index set in any order. Note \(\mu(X) = \sum_{n=1}^\infty 2^{-n} = 1 < \infty\). Define \(f : X \to [0, \infty)\) by \(f(n) = n\); every function on \(X\) is \(\mathcal{S}\)-measurable, and \(f\) is unbounded.

First, \(\int f \, d\mu < \infty\). Let \(f_N = \sum_{n=1}^{N} n \chi_{\{n\}}\). Then \(0 \le f_1 \le f_2 \le \cdots\) and \(\lim_{N \to \infty} f_N(m) = f(m)\) for every \(m \in X\) (indeed \(f_N(m) = m\) once \(N \ge m\)). By the Monotone Convergence Theorem (3.11) and the formula 3.7 for the integral of a simple function,

\begin{equation*} \int f \, d\mu = \lim_{N \to \infty} \int f_N \, d\mu = \lim_{N \to \infty} \sum_{n=1}^{N} n \, \mu(\{n\}) = \sum_{n=1}^{\infty} \frac{n}{2^{n}} = 2 < \infty . \end{equation*}

Now let \(A_1, \ldots, A_m\) be any \(\mathcal{S}\)-partition of \(X\). Since \(X = \mathbb{Z}^+\) is infinite and is the union of the finitely many sets \(A_1, \ldots, A_m\), at least one of them, say \(A_k\), is infinite. Then \(A_k\) contains arbitrarily large positive integers, so \(\sup_{A_k} f = \infty\). Also \(A_k \neq \varnothing\), so picking any \(n_0 \in A_k\) gives \(\mu(A_k) \ge 2^{-n_0} > 0\). Hence \(\mu(A_k) \sup_{A_k} f = \infty\), and therefore

\begin{equation*} \sum_{j=1}^{m} \mu(A_j) \sup_{A_j} f = \infty . \end{equation*}

Because this holds for every \(\mathcal{S}\)-partition, the infimum on the right side of the equation in (a) equals \(\infty\), while \(\int f \, d\mu = 2\). Thus the conclusion of (a) fails, even though \(\mu(X) < \infty\) and \(\int f \, d\mu < \infty\); the hypothesis that \(f\) be bounded cannot be dropped.

Part (c).

Let \(X = \mathbb{Z}^+\), let \(\mathcal{S}\) be the \(\sigma\)-algebra of all subsets of \(\mathbb{Z}^+\), and let \(\mu\) be counting measure, so \(\mu(X) = \infty\). Define \(f : X \to [0, \infty)\) by \(f(n) = 2^{-n}\); this \(f\) is \(\mathcal{S}\)-measurable and bounded, since \(0 < f(n) \le \frac12\) for all \(n\).

As in part (b), put \(f_N = \sum_{n=1}^{N} 2^{-n} \chi_{\{n\}}\). These form an increasing sequence of nonnegative simple functions converging pointwise to \(f\), so the Monotone Convergence Theorem (3.11) together with 3.7 gives

\begin{equation*} \int f \, d\mu = \lim_{N \to \infty} \sum_{n=1}^{N} 2^{-n} \mu(\{n\}) = \sum_{n=1}^{\infty} 2^{-n} = 1 . \end{equation*}

(This is the special case of the assertion in Example 3.6 that integration with respect to counting measure is summation.)

Now let \(A_1, \ldots, A_m\) be any \(\mathcal{S}\)-partition of \(X\). As before, at least one \(A_k\) is infinite, so \(\mu(A_k) = \infty\). Since \(A_k \neq \varnothing\) and \(f > 0\) everywhere, \(\sup_{A_k} f \ge f(n_0) > 0\) for any \(n_0 \in A_k\). Hence \(\mu(A_k) \sup_{A_k} f = \infty\) and so

\begin{equation*} \sum_{j=1}^{m} \mu(A_j) \sup_{A_j} f = \infty \end{equation*}

for every \(\mathcal{S}\)-partition. The infimum on the right side of the equation in (a) is therefore \(\infty\), while \(\int f \, d\mu = 1\). Thus the conclusion of (a) fails, even though \(f\) is bounded; the hypothesis \(\mu(X) < \infty\) cannot be deleted.

Problem (3B.5)

Let \(\lambda\) denote Lebesgue measure on \(\mathbb{R}\). Suppose \(f : \mathbb{R} \to \mathbb{R}\) is a Borel measurable function such that \(\int |f| \, d\lambda < \infty\). Prove that

\begin{equation*} \lim_{k \to \infty} \int_{[-k, \, k]} f \, d\lambda = \int f \, d\lambda . \end{equation*}

Solution

For \(k \in \mathbb{Z}^+\) define \(f_k : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} f_k = \chi_{[-k, \, k]} f , \end{equation*}

so that, by the definition 3.24 of integration over a subset,

\begin{equation*} \int_{[-k, \, k]} f \, d\lambda = \int f_k \, d\lambda . \end{equation*}

Each \(f_k\) is Borel measurable, being the product of the Borel measurable function \(f\) with the characteristic function of the Borel set \([-k, k]\) (see 2.46 for products of measurable functions).

We verify the hypotheses of the Dominated Convergence Theorem (3.31) with dominating function \(g = |f|\).

Pointwise convergence: fix \(x \in \mathbb{R}\). Choose \(K \in \mathbb{Z}^+\) with \(K \ge |x|\). Then for every \(k \ge K\) we have \(x \in [-k, k]\) and hence \(f_k(x) = f(x)\). Thus

\begin{equation*} \lim_{k \to \infty} f_k(x) = f(x) \end{equation*}

for every \(x \in \mathbb{R}\); in fact the sequence is eventually constant at each point.

Domination: for every \(k \in \mathbb{Z}^+\) and every \(x \in \mathbb{R}\),

\begin{equation*} |f_k(x)| = \chi_{[-k, \, k]}(x) \, |f(x)| \le |f(x)| = g(x) , \end{equation*}

since \(\chi_{[-k,k]}(x) \in \{0, 1\}\). The function \(g = |f|\) is Borel measurable and nonnegative, and by hypothesis \(\int g \, d\lambda = \int |f| \, d\lambda < \infty\).

All hypotheses of 3.31 are satisfied, so

\begin{equation*} \lim_{k \to \infty} \int f_k \, d\lambda = \int f \, d\lambda , \end{equation*}

which by the first display is exactly

\begin{equation*} \lim_{k \to \infty} \int_{[-k, \, k]} f \, d\lambda = \int f \, d\lambda , \end{equation*}

as desired. (Note that all the integrals appearing here are defined and finite: \(|f_k| \le |f|\) and \(\int |f| \, d\lambda < \infty\), so by 3.8 each of \(\int f_k^+ \, d\lambda\), \(\int f_k^- \, d\lambda\), \(\int f^+ \, d\lambda\), \(\int f^- \, d\lambda\) is finite, and 3.18 applies.)

Problem (3B.6)

Let \(\lambda\) denote Lebesgue measure on \(\mathbb{R}\). Give an example of a continuous function \(f : [0, \infty) \to \mathbb{R}\) such that \(\lim_{t \to \infty} \int_{[0, \, t]} f \, d\lambda\) exists (in \(\mathbb{R}\)) but \(\int_{[0, \, \infty)} f \, d\lambda\) is not defined.

Solution

Let \(\varphi : \mathbb{R} \to [0, 1]\) be the tent function

\begin{equation*} \varphi(s) = \max\{ 0, \; 1 - |2s - 1| \} , \end{equation*}

so \(\varphi\) is continuous, \(\varphi(s) = 0\) for \(s \notin (0,1)\), \(\varphi(\tfrac12) = 1\), and \(\varphi\) is piecewise linear on \([0,1]\); its graph over \([0,1]\) is a triangle of base \(1\) and height \(1\), so its Riemann integral over \([0,1]\) equals \(\frac12\).

For \(n \in \mathbb{Z}^+\) define \(g_n : [0, \infty) \to \mathbb{R}\) by

\begin{equation*} g_n(x) = \frac{(-1)^{n+1}}{n} \, \varphi(x - n) , \end{equation*}

so \(g_n\) is continuous, \(g_n\) vanishes outside the interval \((n, n+1)\), and \(|g_n| \le \frac1n\) everywhere. Finally define \(f : [0, \infty) \to \mathbb{R}\) by

\begin{equation*} f(x) = \sum_{n=1}^{\infty} g_n(x) . \end{equation*}

For each \(x\) at most one term of this sum is nonzero (the intervals \((n, n+1)\) are disjoint), so \(f\) is well defined and real valued.

\(f\) is continuous. Fix \(N \in \mathbb{Z}^+\). If \(n > N\) then \(g_n\) vanishes on \([0, N+1]\), because \((n, n+1) \cap [0, N+1] = \varnothing\). Hence

\begin{equation*} f(x) = \sum_{n=1}^{N} g_n(x) \quad \text{for all } x \in [0, N+1] , \end{equation*}

and the right side is a finite sum of continuous functions, hence continuous. In particular \(f\) agrees with a continuous function on the set \([0, N+1)\), which is open in \([0, \infty)\). Since continuity is a local property and every point of \([0, \infty)\) lies in \([0, N+1)\) for all large \(N\), the function \(f\) is continuous on \([0, \infty)\). Being continuous, \(f\) is Borel measurable (2.41).

Computing \(\int g_n \, d\lambda\). Since \(g_n\) vanishes outside \([n, n+1]\) and is continuous, hence Riemann integrable on \([n, n+1]\), the identification 3.34 of Riemann and Lebesgue integrals gives

\begin{equation*} \int g_n \, d\lambda = \int_{[n, \, n+1]} g_n \, d\lambda = \frac{(-1)^{n+1}}{n} \int_{0}^{1} \varphi = \frac{(-1)^{n+1}}{2n} . \end{equation*}

In particular \(\int |g_n| \, d\lambda = \frac{1}{2n} < \infty\).

The limit exists. Let \(N \in \mathbb{Z}^+\). On \([0, N]\) we have \(f = \sum_{n=1}^{N-1} g_n\) by the displayed identity above, so additivity of integration (3.21), applicable because each \(\int |g_n| \, d\lambda < \infty\), gives

\begin{equation*} \int_{[0, \, N]} f \, d\lambda = \sum_{n=1}^{N-1} \int_{[0, \, N]} g_n \, d\lambda = \sum_{n=1}^{N-1} \frac{(-1)^{n+1}}{2n} . \end{equation*}

By the alternating series test the series \(\sum_{n=1}^{\infty} \frac{(-1)^{n+1}}{2n}\) converges; its sum is \(\frac{\ln 2}{2}\). Now take any real \(t \ge 1\) and set \(N = \lfloor t \rfloor\). Then \([0, t]\) is the disjoint union of \([0, N]\) and \((N, t]\), and \((N, t] \subseteq (N, N+1]\), where \(|f| \le \frac{1}{N}\). Hence by 3.21 and 3.25,

\begin{equation*} \left| \int_{[0, \, t]} f \, d\lambda - \int_{[0, \, N]} f \, d\lambda \right| = \left| \int_{(N, \, t]} f \, d\lambda \right| \le \lambda\bigl( (N, t] \bigr) \sup_{(N, \, t]} |f| \le 1 \cdot \frac{1}{N} . \end{equation*}

As \(t \to \infty\) we have \(N = \lfloor t \rfloor \to \infty\), so the right side tends to \(0\) and \(\int_{[0, N]} f \, d\lambda\) tends to \(\frac{\ln 2}{2}\). Therefore

\begin{equation*} \lim_{t \to \infty} \int_{[0, \, t]} f \, d\lambda = \frac{\ln 2}{2} \in \mathbb{R} . \end{equation*}

The integral over \([0, \infty)\) is not defined. Note \(g_n \ge 0\) when \(n\) is odd and \(g_n \le 0\) when \(n\) is even, and the \(g_n\) have disjoint supports; hence

\begin{equation*} f^+ = \sum_{n \text{ odd}} g_n \quad \text{and} \quad f^- = - \sum_{n \text{ even}} g_n . \end{equation*}

Fix \(K \in \mathbb{Z}^+\). Then \(f^+ \ge \sum_{k=1}^{K} g_{2k-1} \ge 0\) pointwise, so by 3.8 and additivity of integration for nonnegative functions (3.16),

\begin{equation*} \int f^+ \, d\lambda \ge \sum_{k=1}^{K} \int g_{2k-1} \, d\lambda = \sum_{k=1}^{K} \frac{1}{2(2k-1)} . \end{equation*}

Since \(\sum_{k=1}^{\infty} \frac{1}{2(2k-1)} = \infty\) (its terms are at least \(\frac{1}{4k}\) and the harmonic series diverges), letting \(K \to \infty\) gives \(\int f^+ \, d\lambda = \infty\). The same argument with the even-indexed pieces gives

\begin{equation*} \int f^- \, d\lambda \ge \sum_{k=1}^{K} \frac{1}{2 \cdot 2k} \quad \text{for every } K, \end{equation*}

so \(\int f^- \, d\lambda = \infty\) as well. Because neither \(\int f^+ \, d\lambda\) nor \(\int f^- \, d\lambda\) is finite, the definition 3.18 of the integral of a real-valued function does not apply: \(\int_{[0, \, \infty)} f \, d\lambda\) is not defined.

Thus \(f\) has both required properties. Intuitively, \(f\) is a continuous version of the alternating harmonic series: the truncated integrals converge because of cancellation between consecutive bumps, but the positive and negative parts each have infinite total area.

Problem (3B.7)

Let \(\lambda\) denote Lebesgue measure on \(\mathbb{R}\). Give an example of a continuous function \(f : (0, 1) \to \mathbb{R}\) such that \(\lim_{n \to \infty} \int_{(\frac{1}{n}, \, 1)} f \, d\lambda\) exists (in \(\mathbb{R}\)) but \(\int_{(0, \, 1)} f \, d\lambda\) is not defined.

Solution

The idea is the same as in Exercise 6, but now the alternating bumps march toward \(0\) instead of toward \(\infty\); because the intervals shrink, the bumps must grow tall.

Again let \(\varphi : \mathbb{R} \to [0,1]\) be the tent function \(\varphi(s) = \max\{0, \, 1 - |2s-1|\}\), which is continuous and vanishes outside \((0,1)\).

For \(n \in \mathbb{Z}^+\) let

\begin{equation*} I_n = \Bigl( \frac{1}{n+1}, \, \frac{1}{n} \Bigr) , \qquad \text{so} \qquad \lambda(I_n) = \frac{1}{n} - \frac{1}{n+1} = \frac{1}{n(n+1)} , \end{equation*}

and define \(h_n : (0,1) \to \mathbb{R}\) by

\begin{equation*} h_n(x) = (-1)^{n+1} \, 2(n+1) \, \varphi\bigl( n(n+1)x - n \bigr) . \end{equation*}

Note that \(\varphi(n(n+1)x - n) \neq 0\) exactly when \(0 < n(n+1)x - n < 1\), that is, exactly when \(\frac{1}{n+1} < x < \frac{1}{n}\). Thus \(h_n\) is continuous on \((0,1)\), vanishes outside \(I_n\), and \(|h_n| \le 2(n+1)\). The graph of \(|h_n|\) over \(I_n\) is a triangle with base \(\lambda(I_n) = \frac{1}{n(n+1)}\) and height \(2(n+1)\). Since \(h_n\) vanishes outside \(I_n\), we have \(\int h_n \, d\lambda = \int_{[\frac{1}{n+1}, \, \frac{1}{n}]} h_n \, d\lambda\), and \(h_n\) restricted to the closed interval \([\frac{1}{n+1}, \frac{1}{n}]\) is continuous, hence Riemann integrable there; so the identification 3.34 of the Riemann and Lebesgue integrals gives

\begin{equation*} \int h_n \, d\lambda = (-1)^{n+1} \cdot \frac{1}{2} \cdot \frac{1}{n(n+1)} \cdot 2(n+1) = \frac{(-1)^{n+1}}{n} , \qquad \int |h_n| \, d\lambda = \frac{1}{n} . \end{equation*}

Now define \(f : (0,1) \to \mathbb{R}\) by

\begin{equation*} f(x) = \sum_{n=1}^{\infty} h_n(x) . \end{equation*}

The intervals \(I_1, I_2, \ldots\) are disjoint, so for each \(x \in (0,1)\) at most one term is nonzero and \(f\) is a well-defined real-valued function.

\(f\) is continuous. Fix \(N \in \mathbb{Z}^+\). If \(n > N\) then \(I_n \subseteq (0, \frac{1}{N+1})\), so \(h_n\) vanishes on \([\frac{1}{N+1}, 1)\). Hence

\begin{equation*} f(x) = \sum_{n=1}^{N} h_n(x) \quad \text{for all } x \in \Bigl( \frac{1}{N+1}, \, 1 \Bigr) , \end{equation*}

and the right side is a finite sum of continuous functions. Since the open sets \((\frac{1}{N+1}, 1)\) for \(N \in \mathbb{Z}^+\) cover \((0,1)\), the function \(f\) is continuous on \((0,1)\). Being continuous, \(f\) is Borel measurable (2.41). (Note that \(f\) is unbounded near \(0\), which is permitted.)

The limit exists. Fix \(n \ge 2\). By the displayed identity with \(N = n - 1\), we have \(f = \sum_{k=1}^{n-1} h_k\) on \((\frac{1}{n}, 1)\). Each \(h_k\) with \(k \le n-1\) vanishes outside \(I_k\), and \(I_k \subseteq (\frac{1}{n}, 1)\) because \(\frac{1}{k+1} \ge \frac{1}{n}\); hence \(\int_{(\frac{1}{n}, 1)} h_k \, d\lambda = \int h_k \, d\lambda\). Since \(\int |h_k| \, d\lambda = \frac1k < \infty\), additivity of integration (3.21) applies and gives

\begin{equation*} \int_{(\frac{1}{n}, \, 1)} f \, d\lambda = \sum_{k=1}^{n-1} \int h_k \, d\lambda = \sum_{k=1}^{n-1} \frac{(-1)^{k+1}}{k} . \end{equation*}

By the alternating series test the series \(\sum_{k=1}^{\infty} \frac{(-1)^{k+1}}{k}\) converges, with sum \(\ln 2\). Therefore

\begin{equation*} \lim_{n \to \infty} \int_{(\frac{1}{n}, \, 1)} f \, d\lambda = \ln 2 \in \mathbb{R} . \end{equation*}

The integral over \((0,1)\) is not defined. Because \(h_k \ge 0\) for \(k\) odd, \(h_k \le 0\) for \(k\) even, and the \(h_k\) have disjoint supports,

\begin{equation*} f^+ = \sum_{k \text{ odd}} h_k \quad \text{and} \quad f^- = - \sum_{k \text{ even}} h_k . \end{equation*}

Fix \(K \in \mathbb{Z}^+\). Then \(f^+ \ge \sum_{k=1}^{K} h_{2k-1} \ge 0\) pointwise on \((0,1)\), so 3.8 together with additivity for nonnegative functions (3.16) gives

\begin{equation*} \int_{(0, \, 1)} f^+ \, d\lambda \ge \sum_{k=1}^{K} \int h_{2k-1} \, d\lambda = \sum_{k=1}^{K} \frac{1}{2k-1} . \end{equation*}

Letting \(K \to \infty\) and using the divergence of \(\sum_{k=1}^{\infty} \frac{1}{2k-1}\) (its terms are at least \(\frac{1}{2k}\)), we get \(\int_{(0,1)} f^+ \, d\lambda = \infty\). Similarly

\begin{equation*} \int_{(0, \, 1)} f^- \, d\lambda \ge \sum_{k=1}^{K} \frac{1}{2k} \quad \text{for every } K \in \mathbb{Z}^+ , \end{equation*}

so \(\int_{(0,1)} f^- \, d\lambda = \infty\). Since neither of these two integrals is finite, the definition 3.18 does not assign a value: \(\int_{(0, 1)} f \, d\lambda\) is not defined.

Thus \(f\) has both required properties.

Problem (3B.8)

Verify the assertion in 3.38.

Solution

Recall the setting of the proof of 3.34. Here \(a < b\) and \(f : [a,b] \to \mathbb{R}\) is bounded. For \(n \in \mathbb{Z}^+\) the partition \(P_n\) cuts \([a,b]\) into \(2^n\) closed subintervals \(I_1, \ldots, I_{2^n}\), each of length \(\delta_n = (b-a)/2^n\), and

\begin{equation*} g_n = \sum_{j=1}^{2^n} \Bigl( \inf_{I_j} f \Bigr) \chi_{I_j}, \qquad h_n = \sum_{j=1}^{2^n} \Bigl( \sup_{I_j} f \Bigr) \chi_{I_j}, \end{equation*}

with the modification described in the proof: at each of the finitely many points \(x\) that lie in two of the intervals \(I_1, \ldots, I_{2^n}\), the value \(g_n(x)\) is redefined to be the infimum of \(f\) over the union of the two intervals containing \(x\), and \(h_n(x)\) is redefined to be the supremum of \(f\) over that union. Finally \(f_L(x) = \lim_{n \to \infty} g_n(x)\) and \(f_U(x) = \lim_{n \to \infty} h_n(x)\), these limits existing because \(g_1 \le g_2 \le \cdots\) and \(h_1 \ge h_2 \ge \cdots\) and \(f\) is bounded.

The assertion 3.38 to be verified is

\begin{equation*} \{ x \in [a,b] : f_U(x) \neq f_L(x) \} = \{ x \in [a,b] : f \text{ is not continuous at } x \}. \end{equation*}

A uniform description of \(g_n\) and \(h_n\). For \(x \in [a,b]\) let

\begin{equation*} J_n(x) = \bigcup \{ I_j : 1 \le j \le 2^n \text{ and } x \in I_j \}. \end{equation*}

Thus \(J_n(x)\) is the union of the one or two intervals of \(P_n\) that contain \(x\); in either case \(J_n(x)\) is a closed interval containing \(x\), of length \(\delta_n\) or \(2\delta_n\). With the modification above in force,

\begin{equation*} g_n(x) = \inf_{J_n(x)} f \qquad \text{and} \qquad h_n(x) = \sup_{J_n(x)} f \end{equation*}

for every \(x \in [a,b]\) and every \(n \in \mathbb{Z}^+\). (This is exactly what the modification buys us: at a point lying in two intervals the unmodified formulas would add two terms rather than give a single infimum or supremum.)

Two consequences will be used repeatedly. First, \(x \in J_n(x)\), so

\begin{equation*} g_n(x) \le f(x) \le h_n(x), \end{equation*}

and letting \(n \to \infty\) gives \(f_L(x) \le f(x) \le f_U(x)\) for all \(x \in [a,b]\); in particular \(f_U(x) = f_L(x)\) forces \(f_L(x) = f(x) = f_U(x)\). Second, because \(g_n(x) \le f_L(x)\) and \(f_U(x) \le h_n(x)\) for every \(n\) (monotonicity of the two sequences), we have

\begin{equation*} 0 \le f_U(x) - f_L(x) \le h_n(x) - g_n(x) = \sup_{J_n(x)} f - \inf_{J_n(x)} f \end{equation*}

for every \(n \in \mathbb{Z}^+\).

Continuity at \(x\) implies \(f_L(x) = f_U(x)\). Suppose \(f\) is continuous at \(x\) and let \(\varepsilon > 0\). Choose \(\eta > 0\) such that \(|f(y) - f(x)| < \varepsilon\) for every \(y \in [a,b]\) with \(|y - x| < \eta\). Choose \(n \in \mathbb{Z}^+\) with \(2\delta_n < \eta\). Since \(J_n(x)\) is an interval containing \(x\) of length at most \(2\delta_n < \eta\), every \(y \in J_n(x)\) satisfies \(|y - x| < \eta\) and hence \(|f(y) - f(x)| < \varepsilon\). Therefore

\begin{equation*} \sup_{J_n(x)} f - \inf_{J_n(x)} f \le 2\varepsilon, \end{equation*}

and so \(0 \le f_U(x) - f_L(x) \le 2\varepsilon\) by the displayed inequality above. As \(\varepsilon > 0\) was arbitrary, \(f_U(x) = f_L(x)\).

\(f_L(x) = f_U(x)\) implies continuity at \(x\). Suppose \(f_L(x) = f_U(x)\); as noted, then \(f_L(x) = f(x) = f_U(x)\). Let \(\varepsilon > 0\). Since \(g_n(x) \to f_L(x)\) and \(h_n(x) \to f_U(x)\) and these two limits are equal, there exists \(n \in \mathbb{Z}^+\) with

\begin{equation*} h_n(x) - g_n(x) < \varepsilon. \end{equation*}

Now we claim that \(x\) is an interior point of \(J_n(x)\) relative to \([a,b]\); that is, there exists \(\eta > 0\) such that

\begin{equation*} \{ y \in [a,b] : |y - x| < \eta \} \subseteq J_n(x). \end{equation*}

Indeed, write \(J_n(x) = [c,d]\). If \(c > a\) then \(c\) is a partition point of \(P_n\) that is not an endpoint of \([a,b]\), so \(c\) lies in two of the intervals \(I_1, \ldots, I_{2^n}\); if \(x\) were equal to \(c\), then \(J_n(x)\) would also contain the interval of \(P_n\) lying immediately to the left of \(c\), contradicting \(\min J_n(x) = c\). Hence \(c > a\) forces \(x > c\). Symmetrically \(d < b\) forces \(x < d\). Consequently either \(x > c\) or \(c = a\), and either \(x < d\) or \(d = b\), which is precisely the statement that \(x\) is interior to \([c,d]\) relative to \([a,b]\). So such an \(\eta > 0\) exists.

Now let \(y \in [a,b]\) with \(|y - x| < \eta\). Then \(y \in J_n(x)\), so

\begin{equation*} g_n(x) = \inf_{J_n(x)} f \le f(y) \le \sup_{J_n(x)} f = h_n(x), \end{equation*}

and the same inequalities hold with \(y\) replaced by \(x\). Subtracting,

\begin{equation*} |f(y) - f(x)| \le h_n(x) - g_n(x) < \varepsilon. \end{equation*}

Thus \(f\) is continuous at \(x\).

The two implications together say that \(f_U(x) = f_L(x)\) if and only if \(f\) is continuous at \(x\), which is 3.38.

Problem (3B.9)

Verify the assertion in Example 3.41.

Solution

The assertion of 3.41 is the following. Suppose \((X, \mathcal{S}, \mu)\) is a measure space, \(E_1, \ldots, E_n\) are disjoint subsets of \(X\), and \(a_1, \ldots, a_n\) are distinct nonzero real numbers. Then

\begin{equation*} a_1 \chi_{E_1} + \cdots + a_n \chi_{E_n} \in L^1(\mu) \end{equation*}

if and only if \(E_k \in \mathcal{S}\) and \(\mu(E_k) < \infty\) for all \(k \in \{1, \ldots, n\}\); and in that case

\begin{equation*} \| a_1 \chi_{E_1} + \cdots + a_n \chi_{E_n} \|_1 = |a_1| \mu(E_1) + \cdots + |a_n| \mu(E_n). \end{equation*}

Note that the sets \(E_k\) are not assumed to be in \(\mathcal{S}\); their measurability is part of what is to be proved in one direction.

Throughout, put

\begin{equation*} f = a_1 \chi_{E_1} + \cdots + a_n \chi_{E_n}, \end{equation*}

a function from \(X\) to \(\mathbb{R}\).

The pointwise description of \(f\). Because \(E_1, \ldots, E_n\) are disjoint, each \(x \in X\) lies in at most one of them, so

\begin{equation*} f(x) = \begin{cases} a_k & \text{if } x \in E_k \text{ for some (necessarily unique) } k, \\ 0 & \text{if } x \notin E_1 \cup \cdots \cup E_n. \end{cases} \end{equation*}

Since \(a_1, \ldots, a_n\) are distinct and nonzero, this shows that for each \(k \in \{1, \ldots, n\}\),

\begin{equation*} E_k = f^{-1}(\{a_k\}). \end{equation*}

Indeed, \(x \in E_k\) gives \(f(x) = a_k\); conversely if \(f(x) = a_k\) then \(f(x) \neq 0\), so \(x\) lies in some \(E_j\) and \(a_j = f(x) = a_k\), which forces \(j = k\) because the numbers \(a_1, \ldots, a_n\) are distinct.

Similarly, by disjointness,

\begin{equation*} |f| = |a_1| \chi_{E_1} + \cdots + |a_n| \chi_{E_n}. \end{equation*}

Suppose \(E_k \in \mathcal{S}\) and \(\mu(E_k) < \infty\) for all \(k\). Each characteristic function \(\chi_{E_k}\) is then \(\mathcal{S}\)-measurable, and a finite linear combination of \(\mathcal{S}\)-measurable real-valued functions is \(\mathcal{S}\)-measurable (by 2.46); hence \(f\) is an \(\mathcal{S}\)-measurable function from \(X\) to \(\mathbb{R}\). The sets \(E_1, \ldots, E_n\) are disjoint sets in \(\mathcal{S}\) and \(|a_1|, \ldots, |a_n| \in [0, \infty)\), so 3.7 applied to the displayed formula for \(|f|\) gives

\begin{equation*} \|f\|_1 = \int |f| \, d\mu = |a_1| \mu(E_1) + \cdots + |a_n| \mu(E_n) < \infty, \end{equation*}

the finiteness holding because each \(\mu(E_k) < \infty\). Thus \(f \in L^1(\mu)\), and the asserted formula for \(\|f\|_1\) holds.

Suppose \(f \in L^1(\mu)\). By the definition of \(L^1(\mu)\) (see 3.40), \(f\) is an \(\mathcal{S}\)-measurable function from \(X\) to \(\mathbb{R}\) with \(\|f\|_1 < \infty\). Fix \(k \in \{1, \ldots, n\}\). The singleton \(\{a_k\}\) is a closed subset of \(\mathbb{R}\) and hence a Borel set, so the definition of \(\mathcal{S}\)-measurability (2.35) gives

\begin{equation*} E_k = f^{-1}(\{a_k\}) \in \mathcal{S}. \end{equation*}

Thus all of \(E_1, \ldots, E_n\) belong to \(\mathcal{S}\), and now 3.7 applies exactly as above and yields

\begin{equation*} |a_1| \mu(E_1) + \cdots + |a_n| \mu(E_n) = \int |f| \, d\mu = \|f\|_1 < \infty. \end{equation*}

Every term on the left is in \([0, \infty]\), so every term is finite; since \(a_k \neq 0\) we have \(|a_k| > 0\), and therefore \(\mu(E_k) < \infty\) for each \(k\).

This proves both implications, and the displayed formula for \(\|f\|_1\) was obtained in each case, completing the verification of 3.41.

A remark on the hypotheses. The assumption that the \(a_k\) are nonzero is needed for the conclusion \(\mu(E_k) < \infty\): if some \(a_k\) were \(0\), then \(E_k\) would contribute nothing to \(f\) and could be a nonmeasurable set or a set of infinite measure. The assumption that the \(a_k\) are distinct is needed for the identification \(E_k = f^{-1}(\{a_k\})\): if \(a_1 = a_2\), then \(f^{-1}(\{a_1\}) = E_1 \cup E_2\), and one could take \(E_1\), \(E_2\) nonmeasurable with measurable union.

Problem (3B.10)

(a) Suppose \((X, \mathcal{S}, \mu)\) is a measure space such that \(\mu(X) < \infty\). Suppose \(p, r\) are positive numbers with \(p < r\). Prove that if \(f : X \to [0, \infty)\) is an \(\mathcal{S}\)-measurable function such that \(\int f^r \, d\mu < \infty\), then \(\int f^p \, d\mu < \infty\).

(b) Give an example to show that the result in (a) can be false without the hypothesis that \(\mu(X) < \infty\).

Solution

(a) First note that \(f^p\) and \(f^r\) are \(\mathcal{S}\)-measurable functions from \(X\) to \([0, \infty)\). Indeed, for \(t > 0\) the map \(s \mapsto s^t\) is continuous on \([0, \infty)\), hence Borel measurable by 2.41, and therefore the composition \(f^t = (s \mapsto s^t) \circ f\) is \(\mathcal{S}\)-measurable by 2.44.

The key pointwise estimate is

\begin{equation*} f(x)^p \le 1 + f(x)^r \qquad \text{for every } x \in X. \end{equation*}

To see this, fix \(x \in X\) and consider two cases. If \(f(x) \le 1\), then \(f(x)^p \le 1 \le 1 + f(x)^r\). If \(f(x) > 1\), then because \(p < r\) and \(f(x) > 1\) we have \(f(x)^p \le f(x)^r \le 1 + f(x)^r\). (The inequality \(s^p \le s^r\) for \(s > 1\) holds since \(s^r / s^p = s^{r-p} \ge 1\) when \(s > 1\) and \(r - p > 0\).) This establishes the estimate in both cases.

Now integrate. The function \(1 + f^r\) is a nonnegative \(\mathcal{S}\)-measurable function, so by the order-preserving property of integration (3.8), additivity of integration for nonnegative functions (3.16), and 3.4 applied to \(\chi_X\),

\begin{equation*} \int f^p \, d\mu \le \int (1 + f^r) \, d\mu = \int 1 \, d\mu + \int f^r \, d\mu = \mu(X) + \int f^r \, d\mu. \end{equation*}

By hypothesis \(\mu(X) < \infty\) and \(\int f^r \, d\mu < \infty\), so the right side is finite. Hence \(\int f^p \, d\mu < \infty\), as desired.

(b) Take \(X = (1, \infty)\), let \(\mathcal{S}\) be the collection of Borel subsets of \((1, \infty)\), and let \(\mu = \lambda\) be Lebesgue measure restricted to \(\mathcal{S}\). Then \(\mu(X) = \infty\). Take

\begin{equation*} p = 1, \qquad r = 2, \qquad f(x) = \frac{1}{x} \quad \text{for } x \in (1, \infty), \end{equation*}

which is a continuous, hence Borel measurable, function from \((1, \infty)\) into \([0, \infty)\).

The functions \(f^2 \chi_{(1,k)}\) for \(k = 2, 3, \ldots\) increase pointwise to \(f^2\) on \((1, \infty)\), so the Monotone Convergence Theorem (3.11), together with 3.34 (which identifies the Lebesgue integral of a continuous function on a bounded closed interval with its Riemann integral), gives

\begin{equation*} \int f^2 \, d\mu = \lim_{k \to \infty} \int_1^k \frac{1}{x^2} \, dx = \lim_{k \to \infty} \Bigl( 1 - \frac{1}{k} \Bigr) = 1 < \infty. \end{equation*}

The same argument applied to \(f\) itself gives

\begin{equation*} \int f^1 \, d\mu = \lim_{k \to \infty} \int_1^k \frac{1}{x} \, dx = \lim_{k \to \infty} \ln k = \infty. \end{equation*}

Thus \(\int f^r \, d\mu < \infty\) while \(\int f^p \, d\mu = \infty\), even though \(p < r\). Hence the conclusion of (a) fails without the hypothesis \(\mu(X) < \infty\).

(An even simpler example, with counting measure: let \(X = \mathbb{Z}^+\), let \(\mathcal{S}\) be all subsets of \(\mathbb{Z}^+\), let \(\mu\) be counting measure, take \(p = 1\), \(r = 2\), and \(f(k) = 1/k\). Then \(\int f^2 \, d\mu = \sum_{k=1}^\infty 1/k^2 < \infty\) but \(\int f \, d\mu = \sum_{k=1}^\infty 1/k = \infty\).)

Problem (3B.11)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(f \in L^1(\mu)\). Prove that the set \(\{ x \in X : f(x) \ne 0 \}\) is the countable union of sets with finite \(\mu\)-measure.

Solution

For each \(k \in \mathbb{Z}^+\) define

\begin{equation*} A_k = \Bigl\{ x \in X : |f(x)| > \tfrac{1}{k} \Bigr\}. \end{equation*}

Because \(f\) is \(\mathcal{S}\)-measurable, so is \(|f|\) (by 2.45), and hence

\begin{equation*} A_k = |f|^{-1}\bigl( (\tfrac{1}{k}, \infty) \bigr) \in \mathcal{S} \end{equation*}

for every \(k \in \mathbb{Z}^+\).

The sets \(A_k\) cover the set in question. Indeed, suppose \(x \in X\) and \(f(x) \ne 0\). Then \(|f(x)| > 0\), so there exists \(k \in \mathbb{Z}^+\) with \(1/k < |f(x)|\), which means \(x \in A_k\). Conversely, if \(x \in A_k\) for some \(k\), then \(|f(x)| > 1/k > 0\) and so \(f(x) \ne 0\). Therefore

\begin{equation*} \{ x \in X : f(x) \ne 0 \} = \bigcup_{k=1}^{\infty} A_k, \end{equation*}

a countable union of sets in \(\mathcal{S}\).

It remains to show that \(\mu(A_k) < \infty\) for each \(k\). Fix \(k \in \mathbb{Z}^+\). From the definition of \(A_k\) we have the pointwise inequality

\begin{equation*} \tfrac{1}{k} \chi_{A_k}(x) \le |f(x)| \qquad \text{for every } x \in X; \end{equation*}

this is clear at points of \(A_k\) (where the left side equals \(1/k < |f(x)|\)) and at points outside \(A_k\) (where the left side equals \(0 \le |f(x)|\)). Both sides are nonnegative \(\mathcal{S}\)-measurable functions, so the order-preserving property of integration (3.8) gives

\begin{equation*} \int \tfrac{1}{k} \chi_{A_k} \, d\mu \le \int |f| \, d\mu = \|f\|_1. \end{equation*}

By 3.7 (with the single set \(A_k \in \mathcal{S}\) and the constant \(1/k\)), the left side equals \(\frac{1}{k} \mu(A_k)\). Hence

\begin{equation*} \mu(A_k) \le k \, \|f\|_1. \end{equation*}

Since \(f \in L^1(\mu)\) we have \(\|f\|_1 < \infty\), and therefore \(\mu(A_k) < \infty\).

Thus \(\{ x \in X : f(x) \ne 0 \}\) is the union of the countably many sets \(A_1, A_2, \ldots\), each of which lies in \(\mathcal{S}\) and has finite \(\mu\)-measure, as required.

Two remarks. First, the sets \(A_k\) increase: \(A_1 \subseteq A_2 \subseteq \cdots\). So the union can also be written as an increasing union, or (replacing \(A_k\) by \(A_k \setminus A_{k-1}\)) as a disjoint countable union of sets of finite measure. Second, the inequality \(\mu(A_k) \le k \|f\|_1\) derived above is the Chebyshev-type estimate \(\mu(\{|f| > t\}) \le \|f\|_1 / t\) for \(t > 0\).

Problem (3B.12)

Suppose

\begin{equation*} f_k(x) = \frac{(1 - x)^k \cos x^k}{\sqrt{x}}. \end{equation*}

Prove that \(\lim_{k \to \infty} \int_0^1 f_k = 0\).

Solution

Let \(\lambda\) denote Lebesgue measure on \(\mathbb{R}\). By 3.39, the assertion to be proved is that

\begin{equation*} \lim_{k \to \infty} \int_{(0,1)} f_k \, d\lambda = 0, \end{equation*}

where each \(f_k\) is regarded as a function on \((0,1)\).

Measurability. Each \(f_k\) is continuous on \((0,1)\) (a quotient of continuous functions with nonvanishing denominator), hence Borel measurable by 2.41, hence Lebesgue measurable.

A dominating function. Define \(g : (0,1) \to (0, \infty)\) by

\begin{equation*} g(x) = \frac{1}{\sqrt{x}}. \end{equation*}

For every \(x \in (0,1)\) we have \(0 < 1 - x < 1\), hence \(0 < (1-x)^k \le 1\), and \(|\cos x^k| \le 1\); therefore

\begin{equation*} |f_k(x)| = \frac{(1-x)^k \, |\cos x^k|}{\sqrt{x}} \le \frac{1}{\sqrt{x}} = g(x) \end{equation*}

for every \(k \in \mathbb{Z}^+\) and every \(x \in (0,1)\).

We now check that \(\int_{(0,1)} g \, d\lambda < \infty\). For \(n \in \mathbb{Z}^+\) let \(g_n = g \chi_{(1/n, 1)}\); note \(g_1 = 0\) since \((1,1) = \emptyset\). The intervals \((1/n, 1)\) increase with \(n\) and \(g \ge 0\), so \(g_1 \le g_2 \le \cdots\), and \(g_n(x) \to g(x)\) for every \(x \in (0,1)\), since each \(x \in (0,1)\) satisfies \(x > 1/n\) for all large \(n\). Because \(g\) is continuous on the closed interval \([1/n, 1]\), the Lebesgue and Riemann integrals agree there (3.34), and

\begin{equation*} \int_{(0,1)} g_n \, d\lambda = \int_{1/n}^{1} \frac{dx}{\sqrt{x}} = \Bigl[ 2\sqrt{x} \Bigr]_{1/n}^{1} = 2 - \frac{2}{\sqrt{n}}. \end{equation*}

By the Monotone Convergence Theorem (3.11),

\begin{equation*} \int_{(0,1)} g \, d\lambda = \lim_{n \to \infty} \Bigl( 2 - \frac{2}{\sqrt{n}} \Bigr) = 2 < \infty. \end{equation*}

Pointwise convergence to \(0\). Fix \(x \in (0,1)\). Then \(0 < 1 - x < 1\), so \((1-x)^k \to 0\) as \(k \to \infty\). Since \(|\cos x^k| \le 1\) and \(1/\sqrt{x}\) does not depend on \(k\),

\begin{equation*} |f_k(x)| \le \frac{(1-x)^k}{\sqrt{x}} \longrightarrow 0 \qquad \text{as } k \to \infty. \end{equation*}

Hence \(\lim_{k \to \infty} f_k(x) = 0\) for every \(x \in (0,1)\).

Conclusion. The hypotheses of the Dominated Convergence Theorem (3.31) are satisfied on the measure space \(\bigl( (0,1), \mathcal{B}, \lambda \bigr)\), where \(\mathcal{B}\) denotes the Borel subsets of \((0,1)\): the functions \(f_1, f_2, \ldots\) are measurable, they converge pointwise to the function \(0\), and they are dominated by the measurable function \(g \ge 0\) with \(\int g \, d\lambda = 2 < \infty\). Therefore

\begin{equation*} \lim_{k \to \infty} \int_0^1 f_k = \lim_{k \to \infty} \int_{(0,1)} f_k \, d\lambda = \int_{(0,1)} 0 \, d\lambda = 0, \end{equation*}

as desired.

(Note that the domination also shows \(f_k \in L^1(\lambda)\) for each \(k\), so each of the integrals \(\int_0^1 f_k\) is a well-defined real number.)

Problem (3B.13)

Give an example of a sequence of nonnegative Borel measurable functions \(f_1, f_2, \ldots\) on \([0,1]\) such that both the following conditions hold.

  • \(\lim_{k \to \infty} \int_0^1 f_k = 0\);
  • \(\sup_{k \ge m} f_k(x) = \infty\) for every \(m \in \mathbb{Z}^+\) and every \(x \in [0,1]\).
Solution

The idea is a “typewriter” sequence whose bumps sweep across \([0,1]\) over and over, getting narrower fast enough that the integrals tend to \(0\) but taller without bound, so that at every point infinitely many of the functions are large.

The construction. For \(n \in \mathbb{Z}^+\) let

\begin{equation*} N_n = 1 + 2 + \cdots + n = \frac{n(n+1)}{2}, \end{equation*}

and set \(N_0 = 0\). Since \(N_0 < N_1 < N_2 < \cdots\) and \(N_n \to \infty\), every \(k \in \mathbb{Z}^+\) can be written uniquely as

\begin{equation*} k = N_{n-1} + j \qquad \text{with } n \in \mathbb{Z}^+ \text{ and } j \in \{1, \ldots, n\}. \end{equation*}

For such \(k\) define \(f_k : [0,1] \to [0, \infty)\) by

\begin{equation*} f_k = \sqrt{n} \; \chi_{[\frac{j-1}{n}, \frac{j}{n}]}. \end{equation*}

So the sequence begins \(1 \cdot \chi_{[0,1]}\); then \(\sqrt{2}\,\chi_{[0,1/2]}, \sqrt{2}\,\chi_{[1/2,1]}\); then \(\sqrt{3}\,\chi_{[0,1/3]}, \sqrt{3}\,\chi_{[1/3,2/3]}, \sqrt{3}\,\chi_{[2/3,1]}\); and so on.

Each \(f_k\) is a nonnegative constant times the characteristic function of a closed interval, hence is a nonnegative Borel measurable function on \([0,1]\).

The first condition. Let \(\lambda\) denote Lebesgue measure. If \(k = N_{n-1} + j\) as above, then by 3.7 (or 3.4 together with 3.7),

\begin{equation*} \int_0^1 f_k \, d\lambda = \sqrt{n} \cdot \lambda\Bigl( \Bigl[ \tfrac{j-1}{n}, \tfrac{j}{n} \Bigr] \Bigr) = \sqrt{n} \cdot \frac{1}{n} = \frac{1}{\sqrt{n}}. \end{equation*}

Write \(n(k)\) for the integer \(n\) attached to \(k\) by the decomposition above. For each fixed \(n\), only the finitely many indices \(k \in \{N_{n-1}+1, \ldots, N_n\}\) have \(n(k) = n\); hence only finitely many \(k\) satisfy \(n(k) \le n\), and therefore \(n(k) \to \infty\) as \(k \to \infty\). Consequently

\begin{equation*} \lim_{k \to \infty} \int_0^1 f_k = \lim_{k \to \infty} \frac{1}{\sqrt{n(k)}} = 0. \end{equation*}

The second condition. Fix \(x \in [0,1]\) and \(m \in \mathbb{Z}^+\). Let \(n \in \mathbb{Z}^+\) be arbitrary. The \(n\) intervals

\begin{equation*} \Bigl[ 0, \tfrac{1}{n} \Bigr], \ \Bigl[ \tfrac{1}{n}, \tfrac{2}{n} \Bigr], \ \ldots, \ \Bigl[ \tfrac{n-1}{n}, 1 \Bigr] \end{equation*}

cover \([0,1]\), so there exists \(j \in \{1, \ldots, n\}\) with \(x \in [\frac{j-1}{n}, \frac{j}{n}]\). Put \(k_n = N_{n-1} + j\). Then

\begin{equation*} f_{k_n}(x) = \sqrt{n}. \end{equation*}

Moreover \(k_n > N_{n-1}\), and \(N_{n-1} \to \infty\) as \(n \to \infty\), so \(k_n \ge m\) for all sufficiently large \(n\). Hence for every \(n\) large enough,

\begin{equation*} \sup_{k \ge m} f_k(x) \ge f_{k_n}(x) = \sqrt{n}. \end{equation*}

Letting \(n \to \infty\) gives \(\sup_{k \ge m} f_k(x) = \infty\). Since \(x \in [0,1]\) and \(m \in \mathbb{Z}^+\) were arbitrary, the second condition holds.

Remark. The second condition says \(\limsup_{k \to \infty} f_k(x) = \infty\) for every \(x \in [0,1]\); in particular \(f_1, f_2, \ldots\) converges pointwise at no point of \([0,1]\), even though the integrals tend to \(0\). Note also that \(\int_0^1 |f_k - 0| \, d\lambda = 1/\sqrt{n(k)} \to 0\), so this sequence converges to \(0\) in \(L^1\)-norm while diverging at every point; this shows that convergence in \(L^1\)-norm does not imply pointwise convergence anywhere.

Problem (3B.14)

Let \(\lambda\) denote Lebesgue measure on \(\mathbb{R}\).

(a) Let \(f(x) = 1/\sqrt{x}\). Prove that \(\int_{[0,1]} f \, d\lambda = 2\).

(b) Let \(f(x) = 1/(1 + x^2)\). Prove that \(\int_{\mathbb{R}} f \, d\lambda = \pi\).

(c) Let \(f(x) = (\sin x)/x\). Show that the integral \(\int_{(0,\infty)} f \, d\lambda\) is not defined but \(\lim_{t \to \infty} \int_{(0,t)} f \, d\lambda\) exists in \(\mathbb{R}\).

Solution

(a) Define \(f : [0,1] \to [0, \infty]\) by \(f(x) = 1/\sqrt{x}\) for \(x \in (0,1]\) and \(f(0) = \infty\). (The value assigned at \(0\) is irrelevant to the integral, since \(\lambda(\{0\}) = 0\); taking \(f(0) = \infty\) makes the monotone approximation below converge at every point of \([0,1]\).)

Measurability. Let \(a \in \mathbb{R}\). If \(a \le 0\) then \(f^{-1}((a, \infty]) = [0,1]\). If \(a > 0\) then, since for \(x \in (0,1]\) we have \(1/\sqrt{x} > a\) if and only if \(x < 1/a^2\),

\begin{equation*} f^{-1}\bigl( (a, \infty] \bigr) = [0,1] \cap \Bigl[ 0, \tfrac{1}{a^2} \Bigr), \end{equation*}

which is \([0, 1/a^2)\) if \(a \ge 1\) and \([0,1]\) if \(0 < a < 1\). In every case the set is a Borel set, so \(f\) is Borel measurable by 2.52.

Monotone approximation. For \(n \in \mathbb{Z}^+\) define \(g_n : [0,1] \to [0, \infty)\) by \(g_n = \min\{f, n\}\); explicitly,

\begin{equation*} g_n(x) = \begin{cases} n & \text{if } 0 \le x \le \dfrac{1}{n^2}, \\[2mm] \dfrac{1}{\sqrt{x}} & \text{if } \dfrac{1}{n^2} \le x \le 1. \end{cases} \end{equation*}

The two formulas agree at \(x = 1/n^2\) (both give \(n\)), so \(g_n\) is continuous on \([0,1]\); in particular \(g_n\) is Borel measurable by 2.41. Clearly \(0 \le g_1 \le g_2 \le \cdots\), and \(g_n(x) \to f(x)\) for every \(x \in [0,1]\): if \(x > 0\) then \(g_n(x) = f(x)\) as soon as \(n \ge 1/\sqrt{x}\), while \(g_n(0) = n \to \infty = f(0)\).

Because \(g_n\) is continuous on \([0,1]\), it is Riemann integrable there, and its Riemann and Lebesgue integrals agree by 3.34. Hence

\begin{equation*} \int_{[0,1]} g_n \, d\lambda = \int_0^1 g_n = n \cdot \frac{1}{n^2} + \int_{1/n^2}^{1} \frac{dx}{\sqrt{x}} = \frac{1}{n} + \Bigl[ 2\sqrt{x} \Bigr]_{1/n^2}^{1} = \frac{1}{n} + \Bigl( 2 - \frac{2}{n} \Bigr) = 2 - \frac{1}{n}. \end{equation*}

By the Monotone Convergence Theorem (3.11),

\begin{equation*} \int_{[0,1]} f \, d\lambda = \lim_{n \to \infty} \int_{[0,1]} g_n \, d\lambda = \lim_{n \to \infty} \Bigl( 2 - \frac{1}{n} \Bigr) = 2. \end{equation*}

(b) Here \(f(x) = 1/(1+x^2)\) is continuous on \(\mathbb{R}\), hence Borel measurable by 2.41, and \(f > 0\).

For \(n \in \mathbb{Z}^+\) let \(g_n = f \chi_{[-n, n]}\). Then \(g_n\) is a nonnegative Borel measurable function, \(g_1 \le g_2 \le \cdots\), and \(g_n(x) \to f(x)\) for every \(x \in \mathbb{R}\) (given \(x\), we have \(g_n(x) = f(x)\) once \(n \ge |x|\)).

Since \(f\) is continuous on \([-n, n]\), it is Riemann integrable there and 3.34 gives

\begin{equation*} \int_{\mathbb{R}} g_n \, d\lambda = \int_{[-n,n]} f \, d\lambda = \int_{-n}^{n} \frac{dx}{1+x^2} = \Bigl[ \arctan x \Bigr]_{-n}^{n} = 2 \arctan n, \end{equation*}

where we used that \(\arctan\) is an antiderivative of \(x \mapsto 1/(1+x^2)\) and that \(\arctan\) is odd.

By the Monotone Convergence Theorem (3.11),

\begin{equation*} \int_{\mathbb{R}} f \, d\lambda = \lim_{n \to \infty} 2 \arctan n = 2 \cdot \frac{\pi}{2} = \pi. \end{equation*}

(c) Let \(f : (0, \infty) \to \mathbb{R}\) be defined by \(f(x) = (\sin x)/x\). This function is continuous on \((0,\infty)\), hence Borel measurable by 2.41. The functions \(s \mapsto \max\{s, 0\}\) and \(s \mapsto \max\{-s, 0\}\) are continuous on \(\mathbb{R}\), hence Borel measurable by 2.41, so \(f^+\) and \(f^-\) are nonnegative Borel measurable functions on \((0,\infty)\) by 2.44.

The integral is not defined. By 3.18, \(\int_{(0,\infty)} f \, d\lambda\) is defined only if at least one of \(\int f^+ \, d\lambda\) and \(\int f^- \, d\lambda\) is finite. We show both are infinite.

For \(k \in \{0, 1, 2, \ldots\}\) set

\begin{equation*} A_k = \bigl( 2k\pi, (2k+1)\pi \bigr) \qquad \text{and} \qquad B_k = \bigl( (2k+1)\pi, (2k+2)\pi \bigr). \end{equation*}

These are open intervals contained in \((0,\infty)\), and \(A_0, B_0, A_1, B_1, \ldots\) are pairwise disjoint (consecutive intervals meet only at the excluded endpoints \(\pi, 2\pi, 3\pi, \ldots\)). On \(A_k\) we have \(\sin x > 0\), so \(f > 0\) there and \(f^+ = f\) on \(A_k\); moreover \(0 < x < (2k+1)\pi\) on \(A_k\), whence

\begin{equation*} f^+(x) = \frac{\sin x}{x} \ge \frac{\sin x}{(2k+1)\pi} \qquad \text{for } x \in A_k. \end{equation*}

Since \(x \mapsto \sin x\) is continuous on the closed interval \([2k\pi, (2k+1)\pi]\), it is Riemann integrable there, and 3.34 together with the Fundamental Theorem of Calculus gives \(\int_{[2k\pi, (2k+1)\pi]} \sin x \, d\lambda = \bigl[ -\cos x \bigr]_{2k\pi}^{(2k+1)\pi} = 1 + 1 = 2\); deleting the two endpoints changes nothing, since a set consisting of two points has \(\lambda\)-measure \(0\) and the integrand is bounded. Hence \(\int \chi_{A_k} \sin x \, d\lambda = 2\), and the order-preserving property 3.8 gives

\begin{equation*} \int f^+ \chi_{A_k} \, d\lambda \ge \frac{2}{(2k+1)\pi}. \end{equation*}

Fix \(n \in \mathbb{Z}^+\). Because \(A_0, \ldots, A_{n-1}\) are disjoint and \(f^+ \ge 0\), we have \(\sum_{k=0}^{n-1} f^+ \chi_{A_k} \le f^+\) pointwise on \((0,\infty)\), so by 3.8 and additivity of integration for nonnegative functions (3.16),

\begin{equation*} \int f^+ \, d\lambda \ \ge \ \sum_{k=0}^{n-1} \int f^+ \chi_{A_k} \, d\lambda \ \ge \ \sum_{k=0}^{n-1} \frac{2}{(2k+1)\pi}. \end{equation*}

The series \(\sum_{k=0}^{\infty} \frac{2}{(2k+1)\pi}\) diverges (its terms are at least \(\frac{1}{(k+1)\pi}\), and the harmonic series diverges). Letting \(n \to \infty\) therefore gives \(\int f^+ \, d\lambda = \infty\).

The same argument with \(B_k\) in place of \(A_k\) handles \(f^-\). On \(B_k\) we have \(\sin x < 0\), so \(f < 0\) and \(f^- = -f\) there, and \(0 < x < (2k+2)\pi\) on \(B_k\), whence \(f^-(x) \ge \frac{-\sin x}{(2k+2)\pi}\) on \(B_k\). Since, exactly as above, \(\int \chi_{B_k} (-\sin x) \, d\lambda = \bigl[ \cos x \bigr]_{(2k+1)\pi}^{(2k+2)\pi} = 1 + 1 = 2\), we get

\begin{equation*} \int f^- \chi_{B_k} \, d\lambda \ge \frac{2}{(2k+2)\pi} = \frac{1}{(k+1)\pi}, \end{equation*}

and summing as before gives \(\int f^- \, d\lambda \ge \sum_{k=0}^{n-1} \frac{1}{(k+1)\pi}\) for every \(n\), so \(\int f^- \, d\lambda = \infty\).

Both \(\int f^+ \, d\lambda\) and \(\int f^- \, d\lambda\) equal \(\infty\), so by 3.18 the integral \(\int_{(0,\infty)} f \, d\lambda\) is not defined.

The limit exists. First observe that \(f\) extends to a continuous function \(\tilde f\) on \([0, \infty)\) by setting \(\tilde f(0) = 1\), because \(\lim_{x \to 0} (\sin x)/x = 1\). Fix \(t > 0\). Then \(\tilde f\) is continuous on \([0,t]\), hence Riemann integrable there. We claim

\begin{equation*} \int_{(0,t)} f \, d\lambda = \int_{[0,t]} \tilde f \, d\lambda = \int_0^t \frac{\sin x}{x} \, dx, \end{equation*}

where the right side denotes a Riemann integral. The second equality is the last assertion of 3.34. For the first: \(\tilde f = f\) on \((0,t)\), and \(\tilde f\) is bounded on \([0,t]\) by some \(c \in (0,\infty)\), so

\begin{equation*} \Bigl| \int_{[0,t]} \tilde f \, d\lambda - \int_{(0,t)} f \, d\lambda \Bigr| = \Bigl| \int \tilde f \chi_{\{0,t\}} \, d\lambda \Bigr| \le \int c \, \chi_{\{0,t\}} \, d\lambda = c \, \lambda(\{0,t\}) = 0 \end{equation*}

by 3.8 and 3.7. (Note also that \(\int_{(0,t)} f \, d\lambda\) is defined, indeed finite, because \(|f| \le c\) on \((0,t)\) and \(\lambda((0,t)) = t < \infty\).)

Now suppose \(t \ge 1\). The functions \(x \mapsto 1/x\) and \(x \mapsto -\cos x\) are continuously differentiable on \([1,t]\), so integration by parts for Riemann integrals gives

\begin{equation*} \int_1^t \frac{\sin x}{x} \, dx = \Bigl[ \frac{-\cos x}{x} \Bigr]_1^t - \int_1^t \frac{\cos x}{x^2} \, dx = \cos 1 - \frac{\cos t}{t} - \int_1^t \frac{\cos x}{x^2} \, dx. \end{equation*}

Consider \(h(x) = (\cos x)/x^2\) on \((1, \infty)\); \(h\) is continuous, hence Borel measurable, and \(|h(x)| \le 1/x^2\). Exactly as in part (a) and (b), the Monotone Convergence Theorem applied to the increasing sequence \(x^{-2} \chi_{(1,n)}\) together with 3.34 gives

\begin{equation*} \int_{(1,\infty)} \frac{1}{x^2} \, d\lambda = \lim_{n \to \infty} \int_1^n \frac{dx}{x^2} = \lim_{n \to \infty} \Bigl( 1 - \frac{1}{n} \Bigr) = 1 < \infty. \end{equation*}

Hence \(h \in L^1\bigl( (1,\infty) \bigr)\).

Let \(t_1, t_2, \ldots\) be any sequence in \((1, \infty)\) with \(t_n \to \infty\). Given \(x \in (1,\infty)\) we have \(t_n > x\) for all large \(n\), so \(h \chi_{(1, t_n)} \to h\) pointwise on \((1,\infty)\); and \(|h \chi_{(1,t_n)}(x)| \le 1/x^2\) for all \(n\), a dominating function with finite integral. Also, \(h\) is continuous on \([1, t_n]\), so 3.34 gives \(\int h \chi_{(1,t_n)} \, d\lambda = \int_1^{t_n} \frac{\cos x}{x^2} \, dx\) with the right side a Riemann integral (again the two endpoints form a \(\lambda\)-null set and \(h\) is bounded there). So the Dominated Convergence Theorem (3.31) gives

\begin{equation*} \lim_{n \to \infty} \int_1^{t_n} \frac{\cos x}{x^2} \, dx = \lim_{n \to \infty} \int h \chi_{(1,t_n)} \, d\lambda = \int_{(1,\infty)} h \, d\lambda. \end{equation*}

Since this holds for every such sequence, the limit

\begin{equation*} L = \lim_{t \to \infty} \int_1^t \frac{\cos x}{x^2} \, dx = \int_{(1,\infty)} \frac{\cos x}{x^2} \, d\lambda \end{equation*}

exists and is a real number.

Also \(|\cos t / t| \le 1/t \to 0\) as \(t \to \infty\). For \(t \ge 1\) additivity of the Riemann integral over adjacent intervals gives \(\int_0^t \frac{\sin x}{x} \, dx = \int_0^1 \frac{\sin x}{x} \, dx + \int_1^t \frac{\sin x}{x} \, dx\) (all three integrands being the continuous function \(\tilde f\)). Therefore, letting \(t \to \infty\) in the integration-by-parts identity,

\begin{equation*} \lim_{t \to \infty} \int_{(0,t)} f \, d\lambda = \int_0^1 \frac{\sin x}{x} \, dx + \lim_{t \to \infty} \int_1^t \frac{\sin x}{x} \, dx = \int_0^1 \frac{\sin x}{x} \, dx + \cos 1 - L, \end{equation*}

which is a real number. Thus \(\lim_{t \to \infty} \int_{(0,t)} f \, d\lambda\) exists in \(\mathbb{R}\), even though \(\int_{(0,\infty)} f \, d\lambda\) is not defined.

(For the record, the value of this limit is \(\pi/2\), but the exercise asks only for existence.)

Problem (3B.15)

Prove or give a counterexample: If \(G\) is an open subset of \((0,1)\), then \(\chi_G\) is Riemann integrable on \([0,1]\).

Solution

The statement is false. We produce an open set \(G \subset (0,1)\) for which \(\chi_G\) is not Riemann integrable on \([0,1]\); the point is that \(G\) can be dense in \((0,1)\) and yet have small Lebesgue measure.

Construction of the counterexample. Let \(r_1, r_2, \ldots\) be an enumeration of the rational numbers in \((0,1)\) and set

\begin{equation*} G = (0,1) \cap \bigcup_{k=1}^{\infty} \left( r_k - \frac{1}{2^{k+2}},\ r_k + \frac{1}{2^{k+2}} \right). \end{equation*}

Then \(G\) is an open subset of \((0,1)\), being the intersection of two open sets. Define \(f = \chi_G\), regarded as a function on \([0,1]\); thus \(f \colon [0,1] \to \mathbb{R}\) is bounded (it takes only the values \(0\) and \(1\)), so the lower and upper Riemann integrals \(L(f,[0,1])\) and \(U(f,[0,1])\) of 1.7 make sense.

Two properties of \(G\) drive the argument.

(i) \(G\) is dense in \((0,1)\), because \(G\) contains every rational number in \((0,1)\): indeed \(r_k \in (r_k - 2^{-k-2}, r_k + 2^{-k-2})\) and \(r_k \in (0,1)\).

(ii) \(\lambda(G) \le \tfrac12\), where \(\lambda\) denotes Lebesgue measure; note \(G\) is open, hence Borel, so \(\lambda(G)\) equals the outer measure \(|G|\) (by 2.68 and 2.69). Indeed, an open interval \(I\) is a one-set cover of itself, so the definition 2.2 of outer measure gives \(|I| \le \ell(I)\); combining this with countable subadditivity of outer measure (2.8) and monotonicity (2.5),

\begin{equation*} \lambda(G) \le \sum_{k=1}^{\infty} \left| \left( r_k - \tfrac{1}{2^{k+2}},\ r_k + \tfrac{1}{2^{k+2}} \right) \right| \le \sum_{k=1}^{\infty} \frac{1}{2^{k+1}} = \frac12 . \end{equation*}

The upper Riemann integral equals \(1\). Let \(P\) be the partition \(0 = x_0 < x_1 < \cdots < x_n = 1\) of \([0,1]\). Fix \(j \in \{1,\ldots,n\}\). The open interval \((x_{j-1}, x_j)\) is a nonempty open subset of \([0,1]\), and \((x_{j-1}, x_j) \cap (0,1)\) is still nonempty and open (it is \((x_{j-1},x_j)\) itself unless \(j = 1\) or \(j = n\), in which cases it is \((0,x_1)\) or \((x_{n-1},1)\), again nonempty and open). By (i), \(G\) meets every nonempty open subset of \((0,1)\), so there exists \(y \in G\) with \(x_{j-1} < y < x_j\). Hence

\begin{equation*} \sup_{[x_{j-1},x_j]} f \ge f(y) = 1, \end{equation*}

and since \(f \le 1\) everywhere we get \(\sup_{[x_{j-1},x_j]} f = 1\). Therefore, by the definition 1.3 of the upper Riemann sum,

\begin{equation*} U(f,P,[0,1]) = \sum_{j=1}^{n} (x_j - x_{j-1}) \cdot 1 = 1 . \end{equation*}

As \(P\) was an arbitrary partition, \(U(f,[0,1]) = \inf_P U(f,P,[0,1]) = 1\).

The lower Riemann integral is at most \(\tfrac12\). Again let \(P\) be the partition \(0 = x_0 < x_1 < \cdots < x_n = 1\). Because \(f\) takes only the values \(0\) and \(1\), for each \(j\) we have

\begin{equation*} \inf_{[x_{j-1},x_j]} f = \begin{cases} 1 & \text{if } [x_{j-1},x_j] \subset G,\\ 0 & \text{otherwise.} \end{cases} \end{equation*}

Let \(J = \{ j : [x_{j-1},x_j] \subset G \}\). The intervals \([x_{j-1},x_j]\) for \(j \in J\) have pairwise disjoint interiors and are all contained in \(G\), so by finite additivity and monotonicity of \(\lambda\) (the overlaps are finitely many points, which have measure \(0\)),

\begin{equation*} \sum_{j \in J} (x_j - x_{j-1}) = \lambda\!\left( \bigcup_{j \in J} [x_{j-1},x_j] \right) \le \lambda(G) \le \frac12 . \end{equation*}

Hence, by the definition 1.3 of the lower Riemann sum,

\begin{equation*} L(f,P,[0,1]) = \sum_{j \in J} (x_j - x_{j-1}) \le \frac12 . \end{equation*}

As \(P\) was arbitrary, \(L(f,[0,1]) = \sup_P L(f,P,[0,1]) \le \tfrac12\).

Conclusion. We have \(L(f,[0,1]) \le \tfrac12 < 1 = U(f,[0,1])\). By 1.9, \(f = \chi_G\) is not Riemann integrable on \([0,1]\). This is the desired counterexample.

Remark. Nothing is wrong with \(\chi_G\) from the Lebesgue point of view: \(G\) is an open, hence Borel, set, so \(\chi_G\) is a Borel measurable bounded function and \(\int_{[0,1]} \chi_G \, d\lambda = \lambda(G)\), a perfectly good finite number. The failure above is exactly the phenomenon that motivates Lebesgue integration. Note also that the hypothesis cannot be salvaged by shrinking \(\lambda(G)\): replacing \(2^{-k-2}\) by \(\varepsilon 2^{-k-1}\) in the construction gives, for every \(\varepsilon \in (0,1)\), a dense open \(G \subset (0,1)\) with \(\lambda(G) \le \varepsilon\), for which \(L(\chi_G,[0,1]) \le \varepsilon\) while \(U(\chi_G,[0,1]) = 1\). On the other hand, if \(G\) is a finite union of open intervals then \(\chi_G\) is Riemann integrable, so the statement fails only because open subsets of \((0,1)\) can be complicated.

Problem (3B.16)

Suppose \(f \in \mathcal{L}^1(\mathbb{R})\).

  • (a) For \(t \in \mathbb{R}\), define \(f_t \colon \mathbb{R} \to \mathbb{R}\) by \(f_t(x) = f(x - t)\). Prove that

\begin{equation*} \lim_{t \to 0} \lVert f - f_t \rVert_1 = 0 . \end{equation*}

  • (b) For \(t > 0\), define \(f_t \colon \mathbb{R} \to \mathbb{R}\) by \(f_t(x) = f(tx)\). Prove that

\begin{equation*} \lim_{t \to 1} \lVert f - f_t \rVert_1 = 0 . \end{equation*}

Solution

Throughout, \(\lambda\) denotes Lebesgue measure and \(\mathcal{S}\) denotes the \(\sigma\)-algebra with respect to which \(\mathcal{L}^1(\mathbb{R})\) is being taken, that is, either the Borel subsets of \(\mathbb{R}\) or the Lebesgue measurable subsets of \(\mathbb{R}\) (3.45). The strategy for both parts is the same: the assertion is obvious for continuous functions that vanish outside a bounded set, and by 3.48 such functions are dense in \(\mathcal{L}^1(\mathbb{R})\). The transfer from the dense class to a general \(f\) requires knowing how the \(L^1\)-norm behaves under the relevant change of variable, so we begin with that.

Step 1: the two \(\sigma\)-algebras are invariant under translations and dilations.

Fix \(t \in \mathbb{R}\) and \(c > 0\), and for \(A \subset \mathbb{R}\) write \(t + A = \{t + a : a \in A\}\) and \(cA = \{ca : a \in A\}\).

Write \(|A|\) for the outer measure of an arbitrary set \(A \subset \mathbb{R}\); recall that \(\lambda(A) = |A|\) for every \(A \in \mathcal{S}\).

Translation invariance of outer measure (2.7) says \(|t + A| = |A|\) for every \(A \subset \mathbb{R}\). For dilations: if \(A \subset \bigcup_{k} I_k\) with each \(I_k\) an open interval, then \(cA \subset \bigcup_k cI_k\) and \(cI_k\) is an open interval whose length is \(c\) times the length of \(I_k\); taking the infimum over all such covers in the definition of outer measure (2.2) gives \(|cA| \le c\,|A|\). Applying this inequality with \(cA\) in place of \(A\) and \(1/c\) in place of \(c\) gives \(|A| \le \frac{1}{c}|cA|\), that is, \(c\,|A| \le |cA|\). Hence

\begin{equation*} |cA| = c\,|A| \qquad \text{for all } A \subset \mathbb{R} \text{ and all } c > 0 . \end{equation*}

Now suppose \(A \in \mathcal{S}\). First take \(\mathcal{S}\) to be the Borel \(\sigma\)-algebra. The collection \(\{A \subset \mathbb{R} : t + A \in \mathcal{S}\}\) is a \(\sigma\)-algebra on \(\mathbb{R}\) (it contains \(\mathbb{R}\) because \(t + \mathbb{R} = \mathbb{R}\), and it is closed under complementation and countable unions because \(t + (\mathbb{R} \setminus A) = \mathbb{R} \setminus (t + A)\) and \(t + \bigcup_k A_k = \bigcup_k (t + A_k)\)) that contains every open set, since a translate of an open set is open; hence it contains every Borel set. The same argument, using that \(cU\) is open when \(U\) is open and \(c > 0\), shows \(cA\) is Borel whenever \(A\) is Borel.

Now take \(\mathcal{S}\) to be the \(\sigma\)-algebra of Lebesgue measurable sets. By 2.70, \(A\) Lebesgue measurable means there is a Borel set \(B \subset A\) with \(|A \setminus B| = 0\). Then \(t + B\) is Borel by the previous paragraph, \(t + B \subset t + A\), and

\begin{equation*} (t+A) \setminus (t+B) = t + (A \setminus B), \quad \text{so} \quad \big| (t+A) \setminus (t+B) \big| = |A \setminus B| = 0 \end{equation*}

by 2.7; thus \(t + A\) is Lebesgue measurable by 2.70. Identically, \(cB\) is Borel, \(cB \subset cA\), and \(|(cA) \setminus (cB)| = |c(A \setminus B)| = c\,|A \setminus B| = 0\), so \(cA\) is Lebesgue measurable. In both cases:

\begin{equation*} A \in \mathcal{S} \ \Longrightarrow\ t + A \in \mathcal{S} \text{ and } cA \in \mathcal{S} . \end{equation*}

Step 2: change of variables in the integral.

Claim. Suppose \(h \in \mathcal{L}^1(\mathbb{R})\).

(i) If \(t \in \mathbb{R}\) and \(h^{(t)}(x) = h(x-t)\), then \(h^{(t)} \in \mathcal{L}^1(\mathbb{R})\) and \(\lVert h^{(t)} \rVert_1 = \lVert h \rVert_1\).

(ii) If \(c > 0\) and \(h_{[c]}(x) = h(cx)\), then \(h_{[c]} \in \mathcal{L}^1(\mathbb{R})\) and \(\lVert h_{[c]} \rVert_1 = \frac{1}{c} \lVert h \rVert_1\).

Proof of the claim. Let \(\tau(x) = x - t\), so that \(h^{(t)} = h \circ \tau\) and \(\tau^{-1}(A) = t + A\) for every \(A \subset \mathbb{R}\). For a Borel set \(B \subset \mathbb{R}\) we have \((h \circ \tau)^{-1}(B) = \tau^{-1}\big(h^{-1}(B)\big) = t + h^{-1}(B) \in \mathcal{S}\) by Step 1, so \(h^{(t)}\) is \(\mathcal{S}\)-measurable.

Next let \(u \colon \mathbb{R} \to [0,\infty]\) be any \(\mathcal{S}\)-measurable function and put \(u^{(t)} = u \circ \tau\); the preimage computation just given applies verbatim to \(u\) (for a Borel subset \(B\) of \([0,\infty]\) we again have \((u \circ \tau)^{-1}(B) = t + u^{-1}(B) \in \mathcal{S}\)), so \(u^{(t)}\) is \(\mathcal{S}\)-measurable and \(\int u^{(t)} \, d\lambda\) is defined. If \(A_1, \ldots, A_m\) is an \(\mathcal{S}\)-partition of \(\mathbb{R}\) (3.1), then \(t + A_1, \ldots, t + A_m\) is again an \(\mathcal{S}\)-partition of \(\mathbb{R}\): the sets lie in \(\mathcal{S}\) by Step 1, they are disjoint, and their union is \(t + \mathbb{R} = \mathbb{R}\). Moreover \(\tau\) maps \(t + A_j\) bijectively onto \(A_j\), so

\begin{equation*} \inf_{t + A_j} u^{(t)} = \inf \{ u(\tau(x)) : x \in t + A_j \} = \inf\{u(y) : y \in A_j\} = \inf_{A_j} u , \end{equation*}

while \(\lambda(t + A_j) = \lambda(A_j)\) by Step 1. Hence the lower Lebesgue sums (3.2) satisfy

\begin{equation*} \mathcal{L}\big(u^{(t)}, \{t+A_j\}_{j=1}^m\big) = \sum_{j=1}^m \lambda(t+A_j) \inf_{t+A_j} u^{(t)} = \sum_{j=1}^m \lambda(A_j) \inf_{A_j} u = \mathcal{L}\big(u, \{A_j\}_{j=1}^m\big). \end{equation*}

Since \(\{A_j\} \mapsto \{t + A_j\}\) is a bijection from the set of \(\mathcal{S}\)-partitions of \(\mathbb{R}\) onto itself (its inverse is \(\{A_j\} \mapsto \{-t + A_j\}\)), taking suprema over all \(\mathcal{S}\)-partitions and using the definition 3.3 of the integral of a nonnegative function gives

\begin{equation*} \int u^{(t)} \, d\lambda = \int u \, d\lambda . \end{equation*}

Applying this with \(u = |h|\) and noting \(|h|^{(t)} = |h^{(t)}|\) yields \(\lVert h^{(t)}\rVert_1 = \lVert h \rVert_1 < \infty\), so \(h^{(t)} \in \mathcal{L}^1(\mathbb{R})\). This proves (i). (Applying it instead with \(u = h^+\) and \(u = h^-\), and using \((h^{(t)})^{\pm} = (h^{\pm})^{(t)}\) together with 3.17 and 3.18, also gives \(\int h^{(t)} \, d\lambda = \int h \, d\lambda\), though we do not need this.)

Part (ii) is proved by the identical argument with \(\tau\) replaced by \(\delta(x) = cx\), for which \(\delta^{-1}(A) = \frac{1}{c}A\) and \(\lambda\big(\frac{1}{c}A\big) = \frac{1}{c}\lambda(A)\) by Step 1. The partition \(A_1, \ldots, A_m\) corresponds to the partition \(\frac{1}{c}A_1, \ldots, \frac{1}{c}A_m\), on which \(\inf_{\frac1c A_j} u_{[c]} = \inf_{A_j} u\), so every lower Lebesgue sum gets multiplied by exactly \(\frac1c\); taking suprema gives \(\int u_{[c]}\, d\lambda = \frac{1}{c}\int u \, d\lambda\) for nonnegative \(\mathcal{S}\)-measurable \(u\), and then \(\lVert h_{[c]}\rVert_1 = \frac1c \lVert h \rVert_1\). This completes the proof of the claim.

Step 3: continuous functions vanishing outside a bounded set are bounded and uniformly continuous.

Suppose \(g \colon \mathbb{R} \to \mathbb{R}\) is continuous and \(\{x \in \mathbb{R} : g(x) \ne 0\}\) is bounded, say \(g(x) = 0\) whenever \(|x| \ge M\), where \(M > 0\). Then \(g\) is continuous on the compact interval \([-M-1, M+1]\), hence bounded there, and hence bounded on \(\mathbb{R}\) (it vanishes off \([-M,M]\)); write \(C = \sup_{\mathbb{R}} |g| < \infty\). Also \(g\) is uniformly continuous on \([-M-1,M+1]\). Given \(\eta > 0\), choose \(\delta \in (0,1)\) such that \(|g(x) - g(y)| < \eta\) whenever \(x, y \in [-M-1,M+1]\) with \(|x - y| < \delta\). Now let \(x, y \in \mathbb{R}\) with \(|x-y| < \delta\). If both lie in \([-M-1,M+1]\) we are done. Otherwise one of them, say \(x\), satisfies \(|x| > M+1\); then \(|y| \ge |x| - |x-y| > M + 1 - 1 = M\), so \(g(x) = g(y) = 0\) and \(|g(x)-g(y)| = 0 < \eta\). Thus \(g\) is uniformly continuous on \(\mathbb{R}\).

Note also that such a \(g\) lies in \(\mathcal{L}^1(\mathbb{R})\): it is continuous, hence Borel measurable, and \(|g| \le C \chi_{[-M,M]}\), so \(\lVert g \rVert_1 \le 2CM < \infty\) by 3.8 and 3.4.

Part (a).

By Step 2(i), \(f_t \in \mathcal{L}^1(\mathbb{R})\) for every \(t \in \mathbb{R}\), so \(\lVert f - f_t \rVert_1\) is a finite nonnegative number. Let \(\varepsilon > 0\).

By 3.48 there exists a continuous \(g \colon \mathbb{R} \to \mathbb{R}\) such that \(\{x : g(x) \ne 0\}\) is bounded and

\begin{equation*} \lVert f - g \rVert_1 < \frac{\varepsilon}{3}. \end{equation*}

Fix \(M > 0\) with \(g(x) = 0\) for \(|x| \ge M\), and for \(t \in \mathbb{R}\) put \(g_t(x) = g(x-t)\).

First we control \(\lVert g - g_t \rVert_1\). Suppose \(|t| \le 1\) and \(|x| > M + 1\). Then \(|x| > M\), so \(g(x) = 0\), and \(|x - t| \ge |x| - |t| > M\), so \(g_t(x) = 0\). Thus \(g - g_t\) vanishes outside \([-M-1, M+1]\), and therefore

\begin{equation*} |g - g_t| \le \Big( \sup_{x \in \mathbb{R}} |g(x) - g(x-t)| \Big) \chi_{[-M-1,\,M+1]} . \end{equation*}

By 3.8 and 3.4 this gives

\begin{equation*} \lVert g - g_t \rVert_1 \le (2M+2) \sup_{x \in \mathbb{R}} |g(x) - g(x-t)| . \end{equation*}

By Step 3, \(g\) is uniformly continuous, so there exists \(\delta > 0\), which we may take to satisfy \(\delta \le 1\), such that

\begin{equation*} |g(x) - g(y)| < \frac{\varepsilon}{3(2M+2)} \quad \text{whenever } |x - y| < \delta . \end{equation*}

Consequently, for every \(t\) with \(|t| < \delta\) we have \(|g(x) - g(x-t)| \le \frac{\varepsilon}{3(2M+2)}\) for all \(x\) (the supremum of quantities each less than that bound is at most that bound), and hence

\begin{equation*} \lVert g - g_t \rVert_1 \le (2M+2) \cdot \frac{\varepsilon}{3(2M+2)} = \frac{\varepsilon}{3}. \end{equation*}

Next, \(f_t - g_t\) is precisely the translate \((f-g)_t\) of \(f - g\), so Step 2(i) applied to \(h = f - g \in \mathcal{L}^1(\mathbb{R})\) gives

\begin{equation*} \lVert f_t - g_t \rVert_1 = \lVert f - g \rVert_1 < \frac{\varepsilon}{3} \end{equation*}

for every \(t \in \mathbb{R}\).

Combining the three estimates with the triangle inequality for \(\lVert \cdot \rVert_1\) (3.43), we obtain, for every \(t\) with \(|t| < \delta\),

\begin{equation*} \lVert f - f_t \rVert_1 \le \lVert f - g \rVert_1 + \lVert g - g_t \rVert_1 + \lVert g_t - f_t \rVert_1 < \frac{\varepsilon}{3} + \frac{\varepsilon}{3} + \frac{\varepsilon}{3} = \varepsilon . \end{equation*}

Since \(\varepsilon > 0\) was arbitrary, \(\lim_{t \to 0} \lVert f - f_t \rVert_1 = 0\), as desired.

Part (b).

By Step 2(ii), \(f_t \in \mathcal{L}^1(\mathbb{R})\) for every \(t > 0\). Let \(\varepsilon > 0\).

By 3.48 there exists a continuous \(g \colon \mathbb{R} \to \mathbb{R}\) such that \(\{x : g(x) \ne 0\}\) is bounded and

\begin{equation*} \lVert f - g \rVert_1 < \frac{\varepsilon}{4}. \end{equation*}

Fix \(M > 0\) with \(g(x) = 0\) for \(|x| \ge M\), and for \(t > 0\) put \(g_t(x) = g(tx)\).

First, since \(f_t - g_t\) is the dilate \((f-g)_t\) of \(f - g\), Step 2(ii) applied to \(h = f-g\) gives

\begin{equation*} \lVert f_t - g_t \rVert_1 = \frac{1}{t} \lVert f - g \rVert_1 . \end{equation*}

Hence for every \(t \ge \frac12\),

\begin{equation*} \lVert f_t - g_t \rVert_1 \le 2 \lVert f - g \rVert_1 < \frac{\varepsilon}{2} . \end{equation*}

Second we control \(\lVert g - g_t \rVert_1\) for \(t\) near \(1\). Suppose \(\frac12 \le t \le 2\) and \(|x| > 2M\). Then \(|x| > M\), so \(g(x) = 0\), and \(|tx| \ge \frac12 |x| > M\), so \(g_t(x) = 0\). Thus \(g - g_t\) vanishes outside \([-2M, 2M]\). For \(|x| \le 2M\) we have

\begin{equation*} |x - tx| = |x| \, |t-1| \le 2M |t - 1| . \end{equation*}

By Step 3, \(g\) is uniformly continuous, so there exists \(\eta > 0\) such that

\begin{equation*} |g(u) - g(v)| < \frac{\varepsilon}{16M} \quad \text{whenever } |u - v| < \eta . \end{equation*}

Put \(\rho = \min\{ \tfrac12, \tfrac{\eta}{2M} \} > 0\) and suppose \(|t - 1| < \rho\). Then \(\frac12 < t < \frac32\), and for every \(x\) with \(|x| \le 2M\) we have \(|x - tx| \le 2M|t-1| < \eta\), hence \(|g(x) - g(tx)| < \frac{\varepsilon}{16M}\). Therefore

\begin{equation*} |g - g_t| \le \frac{\varepsilon}{16M} \, \chi_{[-2M,\,2M]} , \end{equation*}

and so by 3.8 and 3.4,

\begin{equation*} \lVert g - g_t \rVert_1 \le 4M \cdot \frac{\varepsilon}{16M} = \frac{\varepsilon}{4}. \end{equation*}

Combining, for every \(t > 0\) with \(|t - 1| < \rho\) (so in particular \(t > \frac12\)), the triangle inequality (3.43) gives

\begin{equation*} \lVert f - f_t \rVert_1 \le \lVert f - g \rVert_1 + \lVert g - g_t \rVert_1 + \lVert g_t - f_t \rVert_1 < \frac{\varepsilon}{4} + \frac{\varepsilon}{4} + \frac{\varepsilon}{2} = \varepsilon . \end{equation*}

Since \(\varepsilon > 0\) was arbitrary, \(\lim_{t \to 1} \lVert f - f_t \rVert_1 = 0\), as desired.

Remark. The dilation constant in Step 2(ii) is what forces the extra care in (b): unlike translation, \(h \mapsto h_{[t]}\) is not an isometry of \(\mathcal{L}^1(\mathbb{R})\), only a bounded operator of norm \(1/t\). Restricting to \(t \ge \frac12\) keeps that factor bounded, which is all that is needed since we are taking \(t \to 1\). Both parts say that the natural actions of the translation group and of the multiplicative group \((0,\infty)\) on \(\mathcal{L}^1(\mathbb{R})\) are strongly continuous; the proofs are the standard \(\varepsilon/3\) argument made possible by 3.48.

Differentiation

Exercises 4A

Problem (4A.1)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(h : X \to \mathbb{R}\) is an \(\mathcal{S}\)-measurable function. Prove that

\begin{equation*} \mu(\{x \in X : |h(x)| \ge c\}) \le \frac{1}{c^p} \int |h|^p \, d\mu \end{equation*}

for all positive numbers \(c\) and \(p\).

Solution

Fix positive numbers \(c\) and \(p\), and set

\begin{equation*} E = \{x \in X : |h(x)| \ge c\}. \end{equation*}

First note that everything in sight is measurable. Because \(h\) is \(\mathcal{S}\)-measurable, so is \(|h|\), and hence so is \(|h|^p\) (the map \(t \mapsto |t|^p\) is a continuous function from \(\mathbb{R}\) to \([0,\infty)\), so it is Borel measurable, and the composition of a Borel measurable function with an \(\mathcal{S}\)-measurable function is \(\mathcal{S}\)-measurable by 2.44 and 2.45). Also

\begin{equation*} E = |h|^{-1}\bigl([c, \infty)\bigr) \in \mathcal{S}, \end{equation*}

because \([c,\infty)\) is a Borel subset of \(\mathbb{R}\).

The point of the proof is the pointwise inequality

\begin{equation*} c^p \chi_E \le |h|^p \quad \text{on all of } X. \end{equation*}

Indeed, if \(x \notin E\) then the left side equals \(0\) and the right side is nonnegative. If \(x \in E\) then \(|h(x)| \ge c > 0\), and because \(t \mapsto t^p\) is increasing on \([0,\infty)\) we get \(|h(x)|^p \ge c^p\), which is exactly what the displayed inequality asserts at \(x\).

Both sides above are nonnegative \(\mathcal{S}\)-measurable functions, so 3.8 (integration is order preserving) applies and gives

\begin{equation*} \int c^p \chi_E \, d\mu \le \int |h|^p \, d\mu . \end{equation*}

The left side equals \(c^p \mu(E)\), since \(c^p \chi_E\) is a simple function whose integral is \(c^p \mu(E)\) by 3.15 (the integral of a linear combination of characteristic functions). Hence

\begin{equation*} c^p \, \mu(\{x \in X : |h(x)| \ge c\}) \le \int |h|^p \, d\mu . \end{equation*}

Dividing by the positive number \(c^p\) gives

\begin{equation*} \mu(\{x \in X : |h(x)| \ge c\}) \le \frac{1}{c^p} \int |h|^p \, d\mu , \end{equation*}

as desired.

Note that no integrability hypothesis was needed: if \(\int |h|^p \, d\mu = \infty\) the inequality is trivially true, and the argument above is valid in that case as well because all the integrals involved are integrals of nonnegative functions and thus are defined in \([0,\infty]\). Taking \(p = 1\) recovers Markov’s inequality (4.1), except that here \(h\) need not belong to \(L^1(\mu)\).

Problem (4A.2)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space with \(\mu(X) = 1\) and \(h \in \mathcal{L}^1(\mu)\). Prove that

\begin{equation*} \mu\left(\left\{x \in X : \left|h(x) - \int h \, d\mu\right| \ge c\right\}\right) \le \frac{1}{c^2}\left(\int h^2 \, d\mu - \left(\int h \, d\mu\right)^2\right) \end{equation*}

for all \(c > 0\).

[The result above is called Chebyshev’s inequality; it plays an important role in probability theory. Pafnuty Chebyshev (1821-1894) was Markov’s thesis advisor.]

Solution

Write

\begin{equation*} m = \int h \, d\mu , \end{equation*}

which is a real number because \(h \in \mathcal{L}^1(\mu)\). Let \(c > 0\).

The function \(h - m\) is a real-valued \(\mathcal{S}\)-measurable function on \(X\). Applying Exercise 1 in this section (with \(p = 2\)) to \(h - m\) gives

\begin{equation*} \mu\bigl(\{x \in X : |h(x) - m| \ge c\}\bigr) \le \frac{1}{c^2}\int (h - m)^2 \, d\mu . \end{equation*}

Thus the proof will be complete once we show

\begin{equation*} \int (h-m)^2 \, d\mu = \int h^2 \, d\mu - m^2 . \end{equation*}

If \(\int h^2 \, d\mu = \infty\), then the right side of the inequality we are proving equals \(\infty\) (because \(m^2 < \infty\)) and there is nothing to prove. So assume from now on that

\begin{equation*} \int h^2 \, d\mu < \infty , \end{equation*}

which says precisely that \(h^2 \in \mathcal{L}^1(\mu)\).

The constant function \(m\) on \(X\) satisfies \(\int |m| \, d\mu = |m| \, \mu(X) = |m| < \infty\), using the hypothesis \(\mu(X) = 1\); thus the constant function \(m\) is in \(\mathcal{L}^1(\mu)\), and so is \(-2mh\) (a scalar multiple of \(h\)). Because

\begin{equation*} (h - m)^2 = h^2 - 2mh + m^2 \end{equation*}

pointwise on \(X\), and each of the three functions on the right side is in \(\mathcal{L}^1(\mu)\), the additivity of integration (3.21), applied twice to handle three summands, together with the homogeneity of integration (3.20), gives

\begin{equation*} \begin{aligned} \int (h-m)^2 \, d\mu &= \int h^2 \, d\mu - 2m \int h \, d\mu + m^2 \, \mu(X) \\ &= \int h^2 \, d\mu - 2m^2 + m^2 \\ &= \int h^2 \, d\mu - m^2 , \end{aligned} \end{equation*}

where we again used \(\mu(X) = 1\) and \(\int h \, d\mu = m\).

Combining the last display with the inequality obtained from Exercise 1 gives

\begin{equation*} \mu\left(\left\{x \in X : \left|h(x) - \int h \, d\mu\right| \ge c\right\}\right) \le \frac{1}{c^2}\left(\int h^2 \, d\mu - \left(\int h \, d\mu\right)^2\right), \end{equation*}

as desired.

In probabilistic language, \(\mu\) is a probability measure, \(h\) is a random variable with expectation \(m\), and the quantity \(\int h^2 \, d\mu - m^2 = \int (h-m)^2 \, d\mu\) is the variance of \(h\); the inequality says that the probability that \(h\) deviates from its mean by at least \(c\) is at most the variance divided by \(c^2\). Note also that the computation above shows the right side is always nonnegative, since it equals the integral of the nonnegative function \((h-m)^2\).

Problem (4A.3)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space. Suppose \(h \in \mathcal{L}^1(\mu)\) and \(\|h\|_1 > 0\). Prove that there is at most one number \(c \in (0, \infty)\) such that

\begin{equation*} \mu(\{x \in X : |h(x)| \ge c\}) = \frac{1}{c}\|h\|_1 . \end{equation*}

Solution

For \(c \in (0,\infty)\) let

\begin{equation*} E_c = \{x \in X : |h(x)| \ge c\} , \end{equation*}

which is in \(\mathcal{S}\) because \(|h|\) is \(\mathcal{S}\)-measurable (2.44 and 2.45). Call \(c\) an equality point if

\begin{equation*} \mu(E_c) = \frac{1}{c}\|h\|_1 . \end{equation*}

Step 1: an equality point forces equality throughout the proof of Markov’s inequality.

Suppose \(c \in (0,\infty)\) is an equality point. Since \(\chi_{E_c}|h| \le |h|\) pointwise and \(|h| = \chi_{E_c}|h| + \chi_{X \setminus E_c}|h|\) pointwise, the additivity of integration for nonnegative functions (3.16) gives

\begin{equation*} \|h\|_1 = \int_{E_c} |h| \, d\mu + \int_{X \setminus E_c} |h| \, d\mu . \end{equation*}

On \(E_c\) we have \(|h| \ge c\), so \(\chi_{E_c}|h| \ge c \chi_{E_c}\) pointwise on \(X\), and hence by 3.8

\begin{equation*} \int_{E_c} |h| \, d\mu \ge c \, \mu(E_c) = \|h\|_1 , \end{equation*}

the last equality being the assumption that \(c\) is an equality point. Substituting this into the previous display yields

\begin{equation*} \|h\|_1 \ge \|h\|_1 + \int_{X \setminus E_c} |h| \, d\mu . \end{equation*}

Because \(\|h\|_1 < \infty\), we may cancel it, obtaining

\begin{equation*} \int_{X \setminus E_c} |h| \, d\mu = 0 \qquad \text{and} \qquad \int_{E_c} |h| \, d\mu = \|h\|_1 = c\,\mu(E_c) . \end{equation*}

Step 2: consequences.

Note first that \(\mu(E_c) = \|h\|_1 / c\) satisfies \(0 < \mu(E_c) < \infty\), because \(0 < \|h\|_1 < \infty\).

Next, the function \(\chi_{E_c}(|h| - c)\) is nonnegative (on \(E_c\) we have \(|h| \ge c\), and off \(E_c\) it is \(0\)) and \(\mathcal{S}\)-measurable, and we claim that its integral is \(0\). Indeed, \(\chi_{E_c}|h| = \chi_{E_c}(|h| - c) + c\chi_{E_c}\) pointwise, so 3.16 gives \(\int_{E_c}|h| \, d\mu = \int \chi_{E_c}(|h| - c) \, d\mu + c\,\mu(E_c)\); since \(\int_{E_c}|h| \, d\mu = c\,\mu(E_c)\) by Step 1 and both of these numbers are finite (as noted in the previous paragraph), the subtraction is legitimate and \(\int \chi_{E_c}(|h| - c) \, d\mu = 0\).

Now we use the standard fact that a nonnegative measurable function with integral \(0\) vanishes almost everywhere. To see this for a nonnegative \(\mathcal{S}\)-measurable \(g\) with \(\int g \, d\mu = 0\): by Exercise 1 in this section applied with \(p = 1\) (or by Markov’s inequality),

\begin{equation*} \mu\bigl(\{x : g(x) \ge \tfrac{1}{n}\}\bigr) \le n \int g \, d\mu = 0 \end{equation*}

for every \(n \in \mathbb{Z}^+\), and \(\{x : g(x) > 0\} = \bigcup_{n=1}^{\infty}\{x : g(x) \ge \frac{1}{n}\}\), so countable subadditivity of \(\mu\) (2.58) gives \(\mu(\{x : g(x) > 0\}) = 0\).

Applying this to \(g = \chi_{E_c}(|h| - c)\) shows

\begin{equation*} |h(x)| = c \quad \text{for almost every } x \in E_c . \end{equation*}

Step 3: at most one equality point.

Suppose \(c_1\) and \(c_2\) are equality points with \(c_1 < c_2\). Because \(c_1 < c_2\) we have \(E_{c_2} \subseteq E_{c_1}\). By Step 2 applied to \(c_1\), the set

\begin{equation*} A = \{x \in E_{c_1} : |h(x)| \ne c_1\} \end{equation*}

has \(\mu(A) = 0\). By Step 2 applied to \(c_2\), the set

\begin{equation*} B = \{x \in E_{c_2} : |h(x)| \ne c_2\} \end{equation*}

has \(\mu(B) = 0\). But every \(x \in E_{c_2}\) lies in \(E_{c_1}\), and \(|h(x)|\) cannot equal both \(c_1\) and \(c_2\) (these are distinct numbers); hence \(E_{c_2} \subseteq A \cup B\). Therefore

\begin{equation*} \mu(E_{c_2}) \le \mu(A) + \mu(B) = 0 \end{equation*}

by subadditivity (2.58). This contradicts \(\mu(E_{c_2}) = \|h\|_1 / c_2 > 0\) from Step 2.

Hence no two distinct numbers in \((0,\infty)\) can both be equality points, so there is at most one \(c \in (0,\infty)\) with \(\mu(\{x \in X : |h(x)| \ge c\}) = \frac{1}{c}\|h\|_1\).

Remark: the hypothesis \(\|h\|_1 > 0\) cannot be dropped. If \(h = 0\), then \(\mu(E_c) = 0 = \frac{1}{c}\|h\|_1\) for every \(c \in (0,\infty)\), so every positive number is an equality point.

Problem (4A.4)

Show that the constant \(3\) in the Vitali Covering Lemma (4.4) cannot be replaced by a smaller positive constant.

Solution

For a bounded nonempty open interval \(I\) and a number \(\alpha > 0\), let \(\alpha * I\) denote the open interval with the same center as \(I\) and \(\alpha\) times the length of \(I\) (this extends Definition 4.2, which is the case \(\alpha = 3\)). We must show: if \(\alpha \in (0,3)\), then there is a list \(I_1, \dots, I_n\) of bounded nonempty open intervals of \(\mathbb{R}\) such that no disjoint sublist \(I_{k_1}, \dots, I_{k_m}\) satisfies

\begin{equation*} I_1 \cup \cdots \cup I_n \subseteq (\alpha * I_{k_1}) \cup \cdots \cup (\alpha * I_{k_m}) . \end{equation*}

So fix \(\alpha\) with \(0 < \alpha < 3\), and choose \(\varepsilon\) with

\begin{equation*} 0 < \varepsilon < \min\{2, \, 3 - \alpha\} . \end{equation*}

Take \(n = 2\) and

\begin{equation*} I_1 = (0, 2), \qquad I_2 = (2 - \varepsilon, 4 - \varepsilon) . \end{equation*}

Then \(I_1\) and \(I_2\) are bounded nonempty open intervals, each of length \(2\), and, because \(0 < \varepsilon < 2\),

\begin{equation*} I_1 \cap I_2 = (2 - \varepsilon, 2) \ne \emptyset , \end{equation*}

so \(I_1, I_2\) is not a disjoint list. Consequently the only disjoint sublists of \(I_1, I_2\) are the sublist consisting of \(I_1\) alone, the sublist consisting of \(I_2\) alone, and the empty sublist.

Also

\begin{equation*} I_1 \cup I_2 = (0, 4 - \varepsilon) . \end{equation*}

The empty sublist clearly fails, because \(I_1 \cup I_2 \ne \emptyset\). Now consider the two one-element sublists.

The interval \(I_1 = (0,2)\) has center \(1\) and length \(2\), so

\begin{equation*} \alpha * I_1 = (1 - \alpha, 1 + \alpha) . \end{equation*}

If \(I_1 \cup I_2 \subseteq \alpha * I_1\), then in particular every point of \((0, 4 - \varepsilon)\) is less than \(1 + \alpha\), so

\begin{equation*} 4 - \varepsilon \le 1 + \alpha , \qquad \text{that is,} \qquad \alpha \ge 3 - \varepsilon . \end{equation*}

But \(\varepsilon < 3 - \alpha\) gives \(\alpha < 3 - \varepsilon\), a contradiction. Hence \(I_1 \cup I_2 \not\subseteq \alpha * I_1\).

The interval \(I_2 = (2 - \varepsilon, 4 - \varepsilon)\) has center \(3 - \varepsilon\) and length \(2\), so

\begin{equation*} \alpha * I_2 = (3 - \varepsilon - \alpha, \; 3 - \varepsilon + \alpha) . \end{equation*}

If \(I_1 \cup I_2 \subseteq \alpha * I_2\), then every point of \((0, 4-\varepsilon)\) exceeds \(3 - \varepsilon - \alpha\), so

\begin{equation*} 3 - \varepsilon - \alpha \le 0 , \qquad \text{that is,} \qquad \alpha \ge 3 - \varepsilon , \end{equation*}

again contradicting \(\alpha < 3 - \varepsilon\). Hence \(I_1 \cup I_2 \not\subseteq \alpha * I_2\).

Thus no disjoint sublist of \(I_1, I_2\) has the covering property, so the conclusion of the Vitali Covering Lemma fails for this list when \(3\) is replaced by \(\alpha\). Since \(\alpha \in (0,3)\) was arbitrary, the constant \(3\) in 4.4 cannot be replaced by any smaller positive constant.

Remark: the example is exactly the extremal configuration of the greedy algorithm used in the proof of 4.4. Two intervals of equal length overlap in an arbitrarily short segment; only one of them can be selected, and stretching it by the factor \(\alpha\) about its center must reach across essentially two full interval lengths beyond its center, i.e. a total half-length of essentially \(3\) times its own half-length. Letting \(\varepsilon \downarrow 0\) shows the constant \(3\) is sharp.

Problem (4A.5)

Prove the assertion left as an exercise in the last sentence of the proof of the Vitali Covering Lemma (4.4).

[The last sentence of that proof asserts: if \(I\) and \(J\) are bounded nonempty open intervals of \(\mathbb{R}\) with \(I \cap J \ne \emptyset\) and \(|I| \ge |J|\), then \(J \subseteq 3 * I\). In the notation of the proof, \(I = I_{k_L}\) and \(J = I_j\).]

Solution

Let \(I\) and \(J\) be bounded nonempty open intervals of \(\mathbb{R}\) with

\begin{equation*} I \cap J \ne \emptyset \qquad \text{and} \qquad |I| \ge |J| , \end{equation*}

where \(|\cdot|\) denotes length. We must show \(J \subseteq 3 * I\).

Write the intervals in terms of their centers and half-lengths:

\begin{equation*} I = (a - r, \, a + r), \qquad J = (b - s, \, b + s) , \end{equation*}

where \(a, b \in \mathbb{R}\) and \(r, s > 0\) (positivity of \(r\) and \(s\) holds because the intervals are nonempty open intervals; boundedness makes \(r\) and \(s\) finite). Then \(|I| = 2r\) and \(|J| = 2s\), so the hypothesis \(|I| \ge |J|\) says

\begin{equation*} s \le r . \end{equation*}

By Definition 4.2,

\begin{equation*} 3 * I = (a - 3r, \, a + 3r) , \end{equation*}

since \(3 * I\) has the same center \(a\) as \(I\) and length \(3|I| = 6r\), hence half-length \(3r\).

First we bound the distance between the two centers. Because \(I \cap J \ne \emptyset\), there exists \(x \in I \cap J\). Then \(|x - a| < r\) and \(|x - b| < s\), so the triangle inequality gives

\begin{equation*} |a - b| \le |a - x| + |x - b| < r + s . \end{equation*}

Now let \(y \in J\) be arbitrary, so \(|y - b| < s\). Then

\begin{equation*} |y - a| \le |y - b| + |b - a| < s + (r + s) = r + 2s \le r + 2r = 3r , \end{equation*}

where the final inequality used \(s \le r\). Thus \(a - 3r < y < a + 3r\), that is, \(y \in 3 * I\).

Since \(y \in J\) was arbitrary, \(J \subseteq 3 * I\), which is the assertion left as an exercise in the proof of 4.4.

Note where each hypothesis was used: the nonempty intersection kept the centers within \(r + s\) of each other, and the inequality \(|I| \ge |J|\) converted the resulting bound \(r + 2s\) into \(3r\). The computation also explains the appearance of the constant \(3\): in the worst case \(s = r\) and the two intervals barely touch, and then \(J\) reaches exactly \(3r\) away from the center of \(I\). (See Exercise 4 in this section, which shows that no constant smaller than \(3\) works.)

Problem (4A.6)

Verify the formula in Example 4.7 for the Hardy-Littlewood maximal function of \(\chi_{[0,1]}\).

[Example 4.7 asserts that, with \(\chi_{[0,1]}\) the characteristic function of the interval \([0,1]\),

\begin{equation*} (\chi_{[0,1]})^*(b) = \begin{cases} \dfrac{1}{2(1-b)} & \text{if } b \le 0, \\[2mm] 1 & \text{if } 0 < b < 1, \\[2mm] \dfrac{1}{2b} & \text{if } b \ge 1. \end{cases} \end{equation*}

as you should verify.]

Solution

Write \(h = \chi_{[0,1]}\). By Definition 4.6,

\begin{equation*} h^*(b) = \sup_{t > 0} \frac{1}{2t}\int_{b-t}^{b+t} |h| = \sup_{t>0} \psi_b(t), \qquad \text{where } \psi_b(t) = \frac{\ell(b,t)}{2t} \end{equation*}

and

\begin{equation*} \ell(b,t) = \bigl|[0,1] \cap [b-t, \, b+t]\bigr| \end{equation*}

is the Lebesgue measure of the overlap (the integral of \(\chi_{[0,1]}\) over \([b-t,b+t]\) is exactly the measure of that intersection). Explicitly, the intersection of the two intervals \([0,1]\) and \([b-t,b+t]\) is \([\max\{b-t,0\}, \min\{b+t,1\}]\) when this is nonempty, so

\begin{equation*} \ell(b,t) = \max\bigl\{0, \; \min\{b+t, 1\} - \max\{b-t, 0\}\bigr\} . \end{equation*}

A useful preliminary observation: since \([0,1] \cap [b-t,b+t] \subseteq [b-t,b+t]\), we always have \(\ell(b,t) \le 2t\), and hence

\begin{equation*} \psi_b(t) \le 1 \quad \text{for all } b \in \mathbb{R} \text{ and } t > 0 . \end{equation*}

Case 1: \(0 < b < 1\).

Let \(t_0 = \min\{b, 1-b\}\), which is positive. If \(0 < t \le t_0\), then \(b - t \ge 0\) and \(b + t \le 1\), so \([b-t,b+t] \subseteq [0,1]\) and therefore \(\ell(b,t) = 2t\), giving \(\psi_b(t) = 1\). Combined with the preliminary observation \(\psi_b \le 1\), this gives

\begin{equation*} h^*(b) = 1 . \end{equation*}

Case 2: \(b \le 0\).

Here \(b - t < b \le 0\) for every \(t > 0\), so \(\max\{b-t,0\} = 0\) and \(\ell(b,t) = \max\{0, \min\{b+t,1\}\}\). We examine three ranges of \(t\).

If \(0 < t \le -b\) (this range is empty when \(b = 0\)), then \(b + t \le 0\), so \(\ell(b,t) = 0\) and \(\psi_b(t) = 0\).

If \(-b < t \le 1 - b\), then \(0 < b + t \le 1\), so \(\ell(b,t) = b + t\) and

\begin{equation*} \psi_b(t) = \frac{b+t}{2t} = \frac{1}{2} + \frac{b}{2t} . \end{equation*}

Because \(b \le 0\), the function \(t \mapsto \frac{b}{2t}\) is nondecreasing on \((0,\infty)\), so \(\psi_b\) is nondecreasing on this range and attains its largest value there at \(t = 1-b\), namely

\begin{equation*} \psi_b(1-b) = \frac{b + (1-b)}{2(1-b)} = \frac{1}{2(1-b)} . \end{equation*}

If \(t > 1 - b\), then \(b + t > 1\), so \(\ell(b,t) = 1\) and \(\psi_b(t) = \frac{1}{2t} < \frac{1}{2(1-b)}\).

Taking the supremum over the three ranges gives

\begin{equation*} h^*(b) = \frac{1}{2(1-b)} , \end{equation*}

and the supremum is attained at \(t = 1 - b\) (note \(1 - b \ge 1 > 0\)).

Case 3: \(b \ge 1\).

Here \(b + t > 1\) for every \(t > 0\), so \(\min\{b+t,1\} = 1\) and \(\ell(b,t) = \max\{0, \, 1 - \max\{b-t,0\}\}\). Again three ranges.

If \(0 < t \le b - 1\) (empty when \(b = 1\)), then \(b - t \ge 1\), so \(\ell(b,t) = 0\) and \(\psi_b(t) = 0\).

If \(b - 1 < t \le b\), then \(0 \le b - t < 1\), so \(\ell(b,t) = 1 - (b-t)\) and

\begin{equation*} \psi_b(t) = \frac{1 - b + t}{2t} = \frac{1}{2} + \frac{1-b}{2t} . \end{equation*}

Because \(1 - b \le 0\), this is nondecreasing in \(t\), so on this range it is largest at \(t = b\), where

\begin{equation*} \psi_b(b) = \frac{1 - b + b}{2b} = \frac{1}{2b} . \end{equation*}

If \(t > b\), then \(b - t < 0\), so \(\ell(b,t) = 1\) and \(\psi_b(t) = \frac{1}{2t} < \frac{1}{2b}\).

Hence

\begin{equation*} h^*(b) = \frac{1}{2b} , \end{equation*}

attained at \(t = b\).

Combining the three cases gives exactly the formula in Example 4.7:

\begin{equation*} (\chi_{[0,1]})^*(b) = \begin{cases} \dfrac{1}{2(1-b)} & \text{if } b \le 0, \\[2mm] 1 & \text{if } 0 < b < 1, \\[2mm] \dfrac{1}{2b} & \text{if } b \ge 1. \end{cases} \end{equation*}

Two remarks. First, the three formulas do not match up at \(b = 0\) and \(b = 1\): the value is \(\frac{1}{2}\) at each of those two points but \(1\) just inside \((0,1)\), so \(h^*\) is discontinuous there. There is no contradiction with Exercise 9 in this section (each set \(\{b : h^*(b) > c\}\) is open): that exercise says only that \(h^*\) is lower semicontinuous, and the jumps at \(b = 0\) and \(b = 1\) are upward as one moves into \((0,1)\), exactly as lower semicontinuity permits. Second, \(h^*\) is not in \(\mathcal{L}^1(\mathbb{R})\), because \(h^*(b) = \frac{1}{2b}\) for \(b \ge 1\) and \(\int_1^{\infty} \frac{db}{2b} = \infty\); this shows that the Hardy-Littlewood maximal inequality (4.8) cannot be strengthened to the assertion that \(h \in \mathcal{L}^1(\mathbb{R})\) implies \(h^* \in \mathcal{L}^1(\mathbb{R})\).

Problem (4A.7)

Find a formula for the Hardy-Littlewood maximal function of the characteristic function of \([0,1] \cup [2,3]\).

Solution

Let \(E = [0,1] \cup [2,3]\) and \(h = \chi_E\). We will show that

\begin{equation*} h^*(b) = \begin{cases} \dfrac{1}{3-b} & \text{if } b \le -1, \\[2mm] \dfrac{1}{2(1-b)} & \text{if } -1 \le b \le 0, \\[2mm] 1 & \text{if } 0 < b < 1, \\[2mm] \dfrac{2b-1}{2b} & \text{if } 1 \le b \le \tfrac{3}{2}, \\[2mm] \dfrac{5-2b}{2(3-b)} & \text{if } \tfrac{3}{2} \le b \le 2, \\[2mm] 1 & \text{if } 2 < b < 3, \\[2mm] \dfrac{1}{2(b-2)} & \text{if } 3 \le b \le 4, \\[2mm] \dfrac{1}{b} & \text{if } b \ge 4. \end{cases} \end{equation*}

(The two formulas agree at each of the overlapping endpoints \(b = -1\), \(b = \frac{3}{2}\), and \(b = 4\), so the formula is unambiguous.)

Setup. By Definition 4.6,

\begin{equation*} h^*(b) = \sup_{t>0} \psi_b(t), \qquad \psi_b(t) = \frac{\bigl|E \cap [b-t, b+t]\bigr|}{2t} . \end{equation*}

Define \(g : \mathbb{R} \to [0,2]\) by \(g(y) = \bigl|E \cap (-\infty, y]\bigr|\), so that

\begin{equation*} g(y) = \begin{cases} 0 & \text{if } y \le 0, \\ y & \text{if } 0 \le y \le 1, \\ 1 & \text{if } 1 \le y \le 2, \\ y - 1 & \text{if } 2 \le y \le 3, \\ 2 & \text{if } y \ge 3. \end{cases} \end{equation*}

Because \(E \cap (-\infty, b+t]\) is the disjoint union of \(E \cap (-\infty, b-t)\) and \(E \cap [b-t, b+t]\), and because \(|E \cap (-\infty, b-t)| = |E \cap (-\infty, b-t]| = g(b-t)\) (a single point has measure \(0\)), additivity of Lebesgue measure gives

\begin{equation*} \bigl|E \cap [b-t, b+t]\bigr| = g(b+t) - g(b-t), \qquad \text{so } \psi_b(t) = \frac{g(b+t) - g(b-t)}{2t} . \end{equation*}

Note that \(g\) is continuous and nondecreasing, and is affine on each of the five intervals \((-\infty,0]\), \([0,1]\), \([1,2]\), \([2,3]\), \([3,\infty)\). Since \(E \cap [b-t,b+t] \subseteq [b-t,b+t]\), we always have

\begin{equation*} \psi_b(t) \le 1 . \end{equation*}

Reduction to finitely many candidates. Fix \(b\) and let

\begin{equation*} T_b = \{t > 0 : b+t \in \{0,1,2,3\} \text{ or } b-t \in \{0,1,2,3\}\} , \end{equation*}

a set with at most eight elements. List its elements as \(t_1 < \cdots < t_N\) and put \(t_0 = 0\), \(t_{N+1} = \infty\). On each interval \([t_i, t_{i+1}] \cap (0,\infty)\), neither \(b+t\) nor \(b-t\) crosses any of the numbers \(0,1,2,3\), so each of \(t \mapsto g(b+t)\) and \(t \mapsto g(b-t)\) is affine there. Hence there are constants \(\alpha, \beta\) (depending on \(i\)) with

\begin{equation*} \psi_b(t) = \frac{\alpha t + \beta}{2t} = \frac{\alpha}{2} + \frac{\beta}{2t} \end{equation*}

on that interval, and this is a monotone function of \(t\). Because \(g\) is continuous, \(\psi_b\) is continuous on \((0,\infty)\), so its supremum over each bounded such interval is attained at an endpoint. Two ends need separate treatment. On \((0, t_1]\) the function \(\psi_b\) is monotone and bounded above by \(1\), so \(\lim_{t \downarrow 0}\psi_b(t)\) exists in \([0,1]\) and the supremum over \((0,t_1]\) equals \(\max\{\lim_{t \downarrow 0}\psi_b(t), \, \psi_b(t_1)\}\). On \([t_N, \infty)\), the bound \(g(b+t) - g(b-t) \le 2\) gives \(\lim_{t \to \infty}\psi_b(t) = 0\); since \(\psi_b \ge 0\), a monotone function on \([t_N,\infty)\) tending to \(0\) is either nonincreasing, in which case its supremum is \(\psi_b(t_N)\), or nondecreasing, in which case it is identically \(0\) there and again its supremum is \(\psi_b(t_N)\). Therefore

\begin{equation*} h^*(b) = \max\Bigl\{\ \lim_{t \downarrow 0}\psi_b(t), \ \psi_b(t_1), \ \dots, \ \psi_b(t_N)\ \Bigr\} . \end{equation*}

Symmetry. The map \(y \mapsto 3 - y\) carries \(E\) onto \(E\) (it interchanges \([0,1]\) and \([2,3]\)) and preserves Lebesgue measure (it is a reflection followed by a translation, each of which preserves the length of every interval and hence preserves outer measure). Applying it to \(E \cap [3-b-t, \, 3-b+t]\) produces \(E \cap [b-t, b+t]\), so \(\psi_{3-b}(t) = \psi_b(t)\) for all \(t > 0\) and therefore

\begin{equation*} h^*(3-b) = h^*(b) \qquad \text{for all } b \in \mathbb{R} . \end{equation*}

Thus it suffices to compute \(h^*(b)\) for \(b \le \frac{3}{2}\); the values for \(b \ge \frac{3}{2}\) follow by reflection.

Case A: \(0 < b < 1\).

For \(0 < t \le \min\{b, 1-b\}\) we have \([b-t,b+t] \subseteq [0,1] \subseteq E\), so \(\psi_b(t) = 1\). Since \(\psi_b \le 1\) always, \(h^*(b) = 1\). By symmetry, \(h^*(b) = 1\) for \(2 < b < 3\) as well.

Case B: \(b \le 0\).

Here \(b - t < 0\) for all \(t > 0\), so \(g(b-t) = 0\) and \(\psi_b(t) = \frac{g(b+t)}{2t}\). The candidate values of \(t\) are those with \(b + t \in \{0,1,2,3\}\), namely \(t = -b\) (only when \(b<0\)), \(t = 1-b\), \(t = 2-b\), \(t = 3-b\); and \(\lim_{t \downarrow 0}\psi_b(t)\) equals \(0\) if \(b < 0\) and equals \(\frac{1}{2}\) if \(b = 0\). Computing:

\begin{equation*} \psi_b(-b) = 0, \quad \psi_b(1-b) = \frac{1}{2(1-b)}, \quad \psi_b(2-b) = \frac{1}{2(2-b)}, \quad \psi_b(3-b) = \frac{2}{2(3-b)} = \frac{1}{3-b} , \end{equation*}

using \(g(1) = g(2) = 1\) and \(g(3) = 2\). Clearly \(\frac{1}{2(2-b)} < \frac{1}{2(1-b)}\), and the limit at \(0\) is at most \(\frac12 = \frac{1}{2(1-b)}\big|_{b=0}\) and in fact never exceeds \(\frac{1}{2(1-b)}\) (it is \(0\) when \(b<0\), and equals \(\frac12\) when \(b = 0\)). Hence

\begin{equation*} h^*(b) = \max\left\{\frac{1}{2(1-b)}, \ \frac{1}{3-b}\right\} . \end{equation*}

Now \(\frac{1}{2(1-b)} \ge \frac{1}{3-b}\) if and only if \(3 - b \ge 2 - 2b\), i.e. if and only if \(b \ge -1\). Therefore

\begin{equation*} h^*(b) = \frac{1}{2(1-b)} \text{ for } -1 \le b \le 0, \qquad h^*(b) = \frac{1}{3-b} \text{ for } b \le -1 . \end{equation*}

Case C: \(1 \le b \le \frac{3}{2}\).

Since \(b \ge 1\), we have \(b + t > 1\) for all \(t > 0\), so the candidates coming from \(b+t\) are \(t = 2-b\) and \(t = 3-b\); those coming from \(b-t\) are \(t = b-1\) (only when \(b > 1\)) and \(t = b\). Also \(\lim_{t \downarrow 0}\psi_b(t)\) equals \(0\) if \(b > 1\) and equals \(\frac{1}{2}\) if \(b = 1\). Using the formula for \(g\) and the ranges \(2b-1 \in [1,2]\), \(2b \in [2,3]\), \(2b-2 \in [0,1]\), \(2b-3 \in [-1,0]\) valid for \(1 \le b \le \frac32\):

\begin{equation*} \begin{aligned} \psi_b(b-1) &= \frac{g(2b-1) - g(1)}{2(b-1)} = \frac{1 - 1}{2(b-1)} = 0 \qquad (b > 1), \\ \psi_b(b) &= \frac{g(2b) - g(0)}{2b} = \frac{(2b-1) - 0}{2b} = \frac{2b-1}{2b}, \\ \psi_b(2-b) &= \frac{g(2) - g(2b-2)}{2(2-b)} = \frac{1 - (2b-2)}{2(2-b)} = \frac{3-2b}{2(2-b)}, \\ \psi_b(3-b) &= \frac{g(3) - g(2b-3)}{2(3-b)} = \frac{2 - 0}{2(3-b)} = \frac{1}{3-b} . \end{aligned} \end{equation*}

Write \(A = \frac{2b-1}{2b} = 1 - \frac{1}{2b}\), \(B = \frac{3-2b}{2(2-b)}\), \(C = \frac{1}{3-b}\).

First, \(A \ge C\): indeed \(A \ge C\) is equivalent to \((2b-1)(3-b) \ge 2b\), i.e. to \(2b^2 - 5b + 3 \le 0\), i.e. to \((2b-3)(b-1) \le 0\), which holds exactly when \(1 \le b \le \frac{3}{2}\).

Second, \(A \ge B\): the function \(b \mapsto A = 1 - \frac{1}{2b}\) is increasing, while \(B = \frac{3-2b}{4-2b}\) has derivative \(\frac{-2(4-2b) + 2(3-2b)}{(4-2b)^2} = \frac{-2}{(4-2b)^2} < 0\), so \(B\) is decreasing; since \(A = B = \frac{1}{2}\) at \(b = 1\), we get \(B \le \frac12 \le A\) on \([1, \frac{3}{2}]\).

Finally the remaining candidates are \(0\) and (when \(b=1\)) \(\frac{1}{2} = A(1)\), neither of which exceeds \(A\). Hence

\begin{equation*} h^*(b) = \frac{2b-1}{2b} \qquad \text{for } 1 \le b \le \tfrac{3}{2} . \end{equation*}

Case D: \(b \ge \frac{3}{2}\), by symmetry.

For \(\frac{3}{2} \le b \le 2\) we have \(1 \le 3 - b \le \frac32\), so by Case C,

\begin{equation*} h^*(b) = h^*(3-b) = \frac{2(3-b) - 1}{2(3-b)} = \frac{5-2b}{2(3-b)} . \end{equation*}

For \(2 < b < 3\) we have \(0 < 3-b < 1\), so \(h^*(b) = 1\) by Case A. For \(3 \le b \le 4\) we have \(-1 \le 3-b \le 0\), so by Case B,

\begin{equation*} h^*(b) = h^*(3-b) = \frac{1}{2\bigl(1 - (3-b)\bigr)} = \frac{1}{2(b-2)} . \end{equation*}

For \(b \ge 4\) we have \(3 - b \le -1\), so by Case B,

\begin{equation*} h^*(b) = h^*(3-b) = \frac{1}{3 - (3-b)} = \frac{1}{b} . \end{equation*}

This establishes the formula displayed at the start.

Sanity checks. The maximum value is \(1\), as it must be since \(0 \le h \le 1\). At \(b = \frac{3}{2}\), the midpoint of the gap, \(h^*(\frac32) = \frac{2}{3}\), achieved with \(t = \frac32\): the interval \((0,3)\) has length \(3\) and meets \(E\) in a set of measure \(2\). At \(b = 1\) and \(b = 2\) the value is \(\frac12\), and at \(b = 0\) and \(b = 3\) it is also \(\frac12\), while \(h^* = 1\) on the open intervals \((0,1)\) and \((2,3)\); thus \(h^*\) is discontinuous exactly at \(0, 1, 2, 3\) and is lower semicontinuous, consistent with Exercise 9 in this section. For large \(|b|\) we get \(h^*(b) \approx \frac{2}{2|b|} = \frac{1}{|b|}\), which is the total measure \(|E| = 2\) divided by the length \(2|b|\) of the smallest centered interval that swallows \(E\); as with Example 4.7, this decay rate shows \(h^* \notin \mathcal{L}^1(\mathbb{R})\).

Problem (4A.8)

Find a formula for the Hardy–Littlewood maximal function of the function \(h \colon \mathbb{R} \to [0, \infty)\) defined by

\begin{equation*} h(x) = \begin{cases} x & \text{if } 0 \le x \le 1,\\ 0 & \text{otherwise.}\end{cases} \end{equation*}

Solution

Because \(h \ge 0\) we have \(|h| = h\). Define

\begin{equation*} F(x) = \int_{-\infty}^{x} h = \begin{cases} 0 & \text{if } x \le 0,\\ \frac{x^2}{2} & \text{if } 0 \le x \le 1,\\ \frac12 & \text{if } x \ge 1,\end{cases} \end{equation*}

so that, in the notation of 4.6,

\begin{equation*} h^{*}(b) = \sup_{t > 0} V_b(t), \qquad V_b(t) := \frac{1}{2t}\int_{b-t}^{b+t} h = \frac{F(b+t) - F(b-t)}{2t}. \end{equation*}

We compute \(\sup_{t>0} V_b(t)\) by splitting the range of \(t\) according to where \(b - t\) and \(b + t\) sit relative to \(0\) and \(1\).

Case 1: \(b \le 0\). Here \(b - t < 0\) for all \(t > 0\), so \(F(b-t) = 0\) and \(V_b(t) = F(b+t)/(2t)\).

If \(t \le -b\) then \(b + t \le 0\) and \(V_b(t) = 0\).

If \(-b \le t \le 1 - b\) then \(b + t \in [0,1]\), so with \(v = b + t \in [0,1]\) (thus \(t = v - b\)) we get

\begin{equation*} V_b(t) = \frac{v^2}{4(v - b)}, \qquad \frac{d}{dv}\Big(\frac{v^2}{v-b}\Big) = \frac{v(v - 2b)}{(v-b)^2} \ge 0 \end{equation*}

because \(v \ge 0\) and \(v - 2b \ge 0\) (as \(b \le 0\)). Hence \(V_b\) increases on this range, with maximum at \(v = 1\), i.e. \(t = 1 - b\), where the value is \(\frac{1}{4(1-b)}\).

If \(t \ge 1 - b\) then \(b + t \ge 1\), so \(V_b(t) = \frac{1/2}{2t} = \frac{1}{4t}\), which is decreasing and again has largest value \(\frac{1}{4(1-b)}\) at \(t = 1-b\).

Therefore \(h^{*}(b) = \frac{1}{4(1-b)}\) for \(b \le 0\).

Case 2: \(0 < b < 1\). For \(0 < t \le \min\{b, 1-b\}\) both endpoints \(b \pm t\) lie in \([0,1]\), so

\begin{equation*} V_b(t) = \frac{(b+t)^2 - (b-t)^2}{4t} = \frac{4bt}{4t} = b . \end{equation*}

So \(V_b \equiv b\) on that initial range. Now we look beyond it.

Suppose first that \(b \le \frac12\), so \(\min\{b,1-b\} = b\). For \(b \le t \le 1 - b\) we have \(b - t \le 0 \le b + t \le 1\), so \(V_b(t) = \frac{(b+t)^2}{4t}\), and

\begin{equation*} \frac{d}{dt}\Big(\frac{(b+t)^2}{4t}\Big) = \frac{(b+t)(t-b)}{4t^2} \ge 0 \quad \text{for } t \ge b, \end{equation*}

so \(V_b\) increases on \([b, 1-b]\), reaching \(\frac{1}{4(1-b)}\) at \(t = 1-b\). For \(t \ge 1-b\) we have \(b+t \ge 1\) and \(b - t \le 0\), so \(V_b(t) = \frac{1}{4t}\), decreasing. Hence for \(0 < b \le \frac12\),

\begin{equation*} h^{*}(b) = \max\Big\{b, \tfrac{1}{4(1-b)}\Big\} = \frac{1}{4(1-b)}, \end{equation*}

the last equality because \(1 - 4b(1-b) = (2b-1)^2 \ge 0\) gives \(\frac{1}{4(1-b)} \ge b\).

Now suppose \(\frac12 \le b < 1\), so \(\min\{b,1-b\} = 1-b\). For \(1 - b \le t \le b\) we have \(b - t \in [0,1]\) and \(b + t \ge 1\), so writing \(u = b - t \in [0, 2b-1]\),

\begin{equation*} V_b(t) = \frac{1 - (b-t)^2}{4t} = \frac{1-u^2}{4(b-u)}, \qquad \frac{d}{du}\Big(\frac{1-u^2}{b-u}\Big) = \frac{u^2 - 2bu + 1}{(b-u)^2}. \end{equation*}

The quadratic \(u^2 - 2bu + 1\) has discriminant \(4b^2 - 4 < 0\), hence is strictly positive; so the expression is strictly increasing in \(u\), i.e. strictly decreasing in \(t\). Thus on \([1-b, b]\) the function \(V_b\) decreases from its value \(b\) at \(t = 1 - b\). For \(t \ge b\) we get \(V_b(t) = \frac{1}{4t} \le \frac{1}{4b} \le b\) (the last step because \(b \ge \frac12\)). Hence \(h^{*}(b) = b\) for \(\frac12 \le b < 1\).

Case 3: \(b \ge 1\). Now \(b + t > 1\) for every \(t>0\), so \(F(b+t) = \frac12\).

If \(t \le b - 1\) then \(b - t \ge 1\) and \(V_b(t) = 0\).

If \(b - 1 \le t \le b\) then \(b - t \in [0,1]\) and, putting \(u = b - t \in [0,1]\),

\begin{equation*} V_b(t) = \frac{1 - u^2}{4(b-u)} =: g(u), \qquad g^{\prime}(u) = \frac{u^2 - 2bu + 1}{4(b-u)^2}. \end{equation*}

The roots of \(u^2 - 2bu + 1\) are \(b \pm \sqrt{b^2-1}\), and \(u_0 := b - \sqrt{b^2-1}\) lies in \((0,1]\) (indeed \(u_0 \le 1\) is equivalent to \(b - 1 \le \sqrt{b^2-1}\), which holds since \((b-1)^2 \le (b-1)(b+1)\)). The other root \(b + \sqrt{b^2-1}\) is at least \(1\), so on \([0,1]\) the quadratic is positive to the left of \(u_0\) and negative to its right; thus \(g^{\prime} > 0\) on \([0,u_0)\) and \(g^{\prime} < 0\) on \((u_0, 1]\), so \(g\) attains its maximum over \([0,1]\) at \(u_0\). Since \(b - u_0 = \sqrt{b^2-1}\) and

\begin{equation*} 1 - u_0^2 = 1 - \big(2b^2 - 1 - 2b\sqrt{b^2-1}\big) = 2\sqrt{b^2-1}\,\big(b - \sqrt{b^2-1}\big), \end{equation*}

we get

\begin{equation*} g(u_0) = \frac{2\sqrt{b^2-1}\,(b-\sqrt{b^2-1})}{4\sqrt{b^2-1}} = \frac{b - \sqrt{b^2-1}}{2}. \end{equation*}

One caveat about the endpoint: the admissible values of \(t\) are \(t > 0\), so \(u = b - t\) ranges over \([0, b) \cap [0,1]\). When \(b > 1\) this is all of \([0,1]\) and the maximum \(g(u_0)\) is attained (at \(t = \sqrt{b^2-1} > 0\)). When \(b = 1\) we have \(u_0 = 1\), which corresponds to the forbidden value \(t = 0\); but \(g\) is increasing on \([0,1)\), so the supremum of \(g\) over \([0,1)\) equals \(\lim_{u \uparrow 1} g(u) = \lim_{u \uparrow 1}\frac{1+u}{4} = \frac12\), which is again the value given by the formula \(\frac{b - \sqrt{b^2-1}}{2}\) at \(b = 1\). So in both cases the supremum over this range of \(t\) equals \(\frac{b-\sqrt{b^2-1}}{2}\) (attained when \(b > 1\), approached as \(t \downarrow 0\) when \(b = 1\)).

If \(t \ge b\) then \(b - t \le 0\) and \(V_b(t) = \frac{1}{4t} \le \frac{1}{4b}\). Since \(b - \sqrt{b^2-1} = \frac{1}{b+\sqrt{b^2-1}} \ge \frac{1}{2b}\), we have \(g(u_0) \ge \frac{1}{4b}\), so this range contributes nothing new. Hence \(h^{*}(b) = \frac{b - \sqrt{b^2-1}}{2}\) for \(b \ge 1\).

Conclusion. Combining the three cases (and noting that at \(b = \frac12\) both formulas give \(\frac12\)):

\begin{equation*} h^{*}(b) = \begin{cases} \dfrac{1}{4(1-b)} & \text{if } b \le \frac12,\\[2mm] b & \text{if } \frac12 \le b < 1,\\[2mm] \dfrac{b - \sqrt{b^2-1}}{2} = \dfrac{1}{2\big(b + \sqrt{b^2-1}\big)} & \text{if } b \ge 1.\end{cases} \end{equation*}

Note that \(h^{*}\) is continuous except at \(b = 1\), where it drops from the left-hand limit \(1\) to the value \(h^{*}(1) = \frac12\); this is consistent with the lower semicontinuity implied by Exercise 9, since the limit from the right at \(1\) also equals \(\frac12\).

Problem (4A.9)

Suppose \(h \colon \mathbb{R} \to \mathbb{R}\) is Lebesgue measurable. Prove that

\begin{equation*} \{b \in \mathbb{R} : h^{*}(b) > c\} \end{equation*}

is an open subset of \(\mathbb{R}\) for every \(c \in \mathbb{R}\).

Solution

Fix \(c \in \mathbb{R}\) and let \(E = \{b \in \mathbb{R} : h^{*}(b) > c\}\). Let \(b \in E\); we produce \(\delta > 0\) with \((b - \delta, b+\delta) \subseteq E\), which shows \(E\) is open.

Because \(h^{*}(b) = \sup_{t>0} \frac{1}{2t}\int_{b-t}^{b+t}|h| > c\), there exists \(t > 0\) such that

\begin{equation*} \frac{1}{2t}\int_{b-t}^{b+t} |h| > c . \end{equation*}

Write \(\alpha = \int_{b-t}^{b+t}|h| \in [0,\infty]\); the displayed inequality says \(\alpha > 2tc\).

Now choose \(t^{\prime} > t\) with \(2t^{\prime}c < \alpha\). Such a \(t^{\prime}\) exists: if \(c \le 0\) then \(2t^{\prime}c \le 0 \le \alpha\) for every \(t^{\prime} > t\), with strict inequality unless \(c = 0\) and \(\alpha = 0\), which cannot happen since \(\alpha > 2tc = 0\) in that case; if \(c > 0\) and \(\alpha = \infty\) take \(t^{\prime} = t+1\); and if \(c > 0\) and \(\alpha < \infty\) then \(\frac{\alpha}{2c} > t\), so any \(t^{\prime}\) in the nonempty interval \(\big(t, \frac{\alpha}{2c}\big)\) works.

Put \(\delta = t^{\prime} - t > 0\) and suppose \(|b^{\prime} - b| < \delta\). Then

\begin{equation*} b^{\prime} - t^{\prime} = b^{\prime} - b - t^{\prime} + b < \delta - t^{\prime} + b = b - t \quad\text{and}\quad b^{\prime} + t^{\prime} = b^{\prime} - b + t^{\prime} + b > -\delta + t^{\prime} + b = b + t, \end{equation*}

so \((b - t, b+t) \subseteq (b^{\prime} - t^{\prime}, b^{\prime} + t^{\prime})\). Since \(|h| \ge 0\), monotonicity of the integral gives

\begin{equation*} \int_{b^{\prime}-t^{\prime}}^{b^{\prime}+t^{\prime}} |h| \ \ge\ \int_{b-t}^{b+t} |h| \ = \ \alpha . \end{equation*}

Therefore

\begin{equation*} h^{*}(b^{\prime}) \ \ge\ \frac{1}{2t^{\prime}}\int_{b^{\prime}-t^{\prime}}^{b^{\prime}+t^{\prime}}|h| \ \ge\ \frac{\alpha}{2t^{\prime}} \ > \ c, \end{equation*}

the last inequality because \(2t^{\prime}c < \alpha\). Hence \(b^{\prime} \in E\).

Thus every point of \(E\) is an interior point, so \(E\) is open. In particular, taking \(c\) to range over \(\mathbb{R}\) shows that \(h^{*}\) is a Borel measurable function, as asserted in the text following 4.7.

Problem (4A.10)

Prove or give a counterexample: If \(h \colon \mathbb{R} \to [0,\infty)\) is an increasing function, then \(h^{*}\) is an increasing function.

Solution

The statement is true; here is a proof.

First note that an increasing function \(h \colon \mathbb{R} \to [0,\infty)\) is Borel measurable (for each \(c \in \mathbb{R}\) the set \(\{x : h(x) > c\}\) is an interval, since \(x \in \{h > c\}\) and \(y > x\) imply \(h(y) \ge h(x) > c\)), so \(h^{*}\) is defined. Also \(|h| = h\) because \(h \ge 0\).

Suppose \(b_1 < b_2\) and let \(s = b_2 - b_1 > 0\). Fix \(t > 0\). Translation invariance of Lebesgue measure (2.7) gives translation invariance of the integral: for every nonnegative Borel measurable \(g\),

\begin{equation*} \int_{b_1 - t}^{b_1+t} g(x + s)\,dx = \int_{b_2-t}^{b_2+t} g(y)\,dy \end{equation*}

(substitute \(y = x+s\); formally, \(x \mapsto g(x+s)\) is the composition of \(g\) with a translation, and the two integrals agree first for characteristic functions of Borel sets by 2.7, hence for simple functions by linearity, hence for all nonnegative Borel measurable functions by the definition of the integral as a supremum over simple functions below the integrand).

Since \(h\) is increasing and \(s > 0\), we have \(h(x) \le h(x+s)\) for every \(x \in \mathbb{R}\). Hence, by monotonicity of the integral and the identity above,

\begin{equation*} \int_{b_1-t}^{b_1+t} h(x)\,dx \ \le\ \int_{b_1-t}^{b_1+t} h(x+s)\,dx \ = \ \int_{b_2-t}^{b_2+t} h(y)\,dy . \end{equation*}

Dividing by \(2t\),

\begin{equation*} \frac{1}{2t}\int_{b_1-t}^{b_1+t} |h| \ \le\ \frac{1}{2t}\int_{b_2-t}^{b_2+t}|h| \ \le\ h^{*}(b_2). \end{equation*}

This holds for every \(t > 0\), so taking the supremum over \(t > 0\) on the left gives \(h^{*}(b_1) \le h^{*}(b_2)\).

Therefore \(h^{*}\) is increasing.

Problem (4A.11)

Give an example of a Borel measurable function \(h \colon \mathbb{R} \to [0,\infty)\) such that \(h^{*}(b) < \infty\) for all \(b \in \mathbb{R}\) but \(\sup\{h^{*}(b) : b \in \mathbb{R}\} = \infty\).

Solution

Define

\begin{equation*} h = \sum_{n=1}^{\infty} n\, \chi_{[n,\, n + n^{-3}]}. \end{equation*}

Because \(n + n^{-3} < n+1\) for every positive integer \(n\), the intervals \([n, n+n^{-3}]\) are pairwise disjoint, so at each \(x \in \mathbb{R}\) at most one term of the sum is nonzero; thus \(h(x) \in [0,\infty)\) for every \(x\). Indeed at each fixed \(x\) the partial sums are eventually constant, so they converge pointwise to \(h\); each partial sum is a finite linear combination of characteristic functions of Borel sets and hence is a Borel measurable real-valued function. Therefore \(h\) is Borel measurable by 2.48.

Step 1: \(h \in \mathcal{L}^1(\mathbb{R})\) and \(h\) is bounded on bounded sets. By the Monotone Convergence Theorem (3.11) applied to the partial sums,

\begin{equation*} \int_{-\infty}^{\infty} h = \sum_{n=1}^{\infty} n \cdot |[n, n+n^{-3}]| = \sum_{n=1}^{\infty} \frac{n}{n^{3}} = \sum_{n=1}^{\infty}\frac{1}{n^{2}} < \infty . \end{equation*}

Write \(\|h\|_1 = \int |h| < \infty\). Moreover, if \(R > 0\) then \(h(x) = 0\) for \(x < 1\) and \(h(x) \le \lceil R \rceil\) for \(x \in [-R, R]\), since the only spikes meeting \([-R,R]\) are those with \(n \le R\). So \(M_R := \sup_{[-R,R]} h < \infty\) for every \(R > 0\).

Step 2: \(h^{*}(b) < \infty\) for every \(b \in \mathbb{R}\). Fix \(b \in \mathbb{R}\) and put \(R = |b| + 1\). If \(0 < t \le 1\) then \([b-t, b+t] \subseteq [-R, R]\), so

\begin{equation*} \frac{1}{2t}\int_{b-t}^{b+t} h \ \le\ \frac{1}{2t}\cdot 2t \cdot M_R = M_R . \end{equation*}

If \(t \ge 1\) then

\begin{equation*} \frac{1}{2t}\int_{b-t}^{b+t} h \ \le\ \frac{1}{2t}\,\|h\|_1 \ \le\ \frac{\|h\|_1}{2}. \end{equation*}

Taking the supremum over all \(t > 0\) gives \(h^{*}(b) \le \max\{M_R, \tfrac12\|h\|_1\} < \infty\).

Step 3: \(\sup_{b \in \mathbb{R}} h^{*}(b) = \infty\). Let \(n\) be a positive integer and take \(b = n\), \(t = n^{-3}\). Then \([n, n+n^{-3}] \subseteq [b-t, b+t]\), so

\begin{equation*} h^{*}(n) \ \ge\ \frac{1}{2n^{-3}}\int_{n - n^{-3}}^{n + n^{-3}} h \ \ge\ \frac{n^{3}}{2}\cdot n \cdot n^{-3} \ = \ \frac{n}{2}. \end{equation*}

Since \(n/2 \to \infty\), we conclude \(\sup\{h^{*}(b) : b \in \mathbb{R}\} = \infty\).

Thus \(h\) has all the required properties.

Problem (4A.12)

Show that \(|\{b \in \mathbb{R} : h^{*}(b) = \infty\}| = 0\) for every \(h \in \mathcal{L}^1(\mathbb{R})\).

Solution

Let \(h \in \mathcal{L}^1(\mathbb{R})\) and set

\begin{equation*} E = \{b \in \mathbb{R} : h^{*}(b) = \infty\}. \end{equation*}

Since \(h^{*}\) takes values in \([0,\infty]\), we have

\begin{equation*} E = \bigcap_{n=1}^{\infty}\{b \in \mathbb{R} : h^{*}(b) > n\}, \end{equation*}

and each set on the right is open by Exercise 9 in this section; hence \(E\) is a Borel set, in particular Lebesgue measurable. (Even without this, the argument below works verbatim with \(|\cdot|\) read as outer measure, which is monotone by 2.5.)

For each positive integer \(n\) we have \(E \subseteq \{b \in \mathbb{R} : h^{*}(b) > n\}\), so the Hardy–Littlewood maximal inequality (4.8) gives

\begin{equation*} |E| \ \le\ \big|\{b \in \mathbb{R} : h^{*}(b) > n\}\big| \ \le\ \frac{3}{n}\,\|h\|_1 . \end{equation*}

Because \(h \in \mathcal{L}^1(\mathbb{R})\), the number \(\|h\|_1\) is finite, so \(\frac{3}{n}\|h\|_1 \to 0\) as \(n \to \infty\). As \(|E|\) is a fixed number in \([0,\infty]\) bounded above by every term of a sequence tending to \(0\), we conclude \(|E| = 0\).

Problem (4A.13)

Show that there exists \(h \in \mathcal{L}^1(\mathbb{R})\) such that \(h^{*}(b) = \infty\) for every \(b \in \mathbb{Q}\).

Solution

Two elementary properties of \(h \mapsto h^{*}\). Let \(f, g \colon \mathbb{R} \to [0,\infty]\) be Borel measurable.

(i) If \(f \le g\) pointwise then \(f^{*} \le g^{*}\); indeed \(\frac{1}{2t}\int_{b-t}^{b+t} f \le \frac{1}{2t}\int_{b-t}^{b+t} g \le g^{*}(b)\) for every \(t > 0\), and we take the supremum over \(t\).

(ii) For \(r \in \mathbb{R}\) let \(g_r(x) = g(x - r)\). Then \(g_r^{*}(b) = g^{*}(b - r)\) for every \(b\), because translation invariance of Lebesgue measure (2.7) gives \(\int_{b-t}^{b+t} g_r = \int_{b-r-t}^{b-r+t} g\) for every \(t > 0\).

A single bad point. Define \(g \colon \mathbb{R} \to [0,\infty)\) by

\begin{equation*} g(x) = \begin{cases} \dfrac{1}{x\,(\log x)^{2}} & \text{if } 0 < x < \frac12,\\[2mm] 0 & \text{otherwise.}\end{cases} \end{equation*}

Then \(g\) is Borel measurable (it is continuous on the open set \((0,\tfrac12)\) and zero elsewhere) and nonnegative. Since \(\frac{d}{dx}\big(\!-\!\frac{1}{\log x}\big) = \frac{1}{x(\log x)^{2}}\) for \(0 < x < 1\), and \(-\frac{1}{\log x} \to 0\) as \(x \downarrow 0\), we get for \(0 < t < \frac12\)

\begin{equation*} \int_{0}^{t} g = -\frac{1}{\log t} = \frac{1}{|\log t|}, \end{equation*}

and in particular \(\|g\|_1 = \int_{0}^{1/2} g = \frac{1}{\log 2} < \infty\), so \(g \in \mathcal{L}^1(\mathbb{R})\). Now for \(0 < t < \frac12\),

\begin{equation*} \frac{1}{2t}\int_{-t}^{t} g = \frac{1}{2t}\int_{0}^{t} g = \frac{1}{2t\,|\log t|} \xrightarrow[\ t \downarrow 0\ ]{} \infty , \end{equation*}

because \(t|\log t| \to 0\) as \(t \downarrow 0\). Hence \(g^{*}(0) = \infty\).

Spreading it over \(\mathbb{Q}\). Let \(r_1, r_2, r_3, \dots\) be an enumeration of \(\mathbb{Q}\) (possible since \(\mathbb{Q}\) is countable), and define \(\tilde h \colon \mathbb{R} \to [0,\infty]\) by

\begin{equation*} \tilde h(x) = \sum_{k=1}^{\infty} 2^{-k} g(x - r_k). \end{equation*}

Each summand \(x \mapsto 2^{-k} g(x - r_k)\) is Borel measurable, because for \(c \in \mathbb{R}\) the set \(\{x : 2^{-k}g(x - r_k) > c\}\) is the translate by \(r_k\) of the Borel set \(\{y : g(y) > 2^{k}c\}\), and translates of Borel sets are Borel. The partial sums \(s_m = \sum_{k=1}^{m} 2^{-k}g(\cdot - r_k)\) are therefore Borel measurable, and since all terms are nonnegative we have \(\tilde h = \sup_m s_m\) pointwise; hence \(\tilde h\) is Borel measurable by 2.53. By the Monotone Convergence Theorem (3.11) applied to the increasing sequence \(s_1 \le s_2 \le \cdots\),

\begin{equation*} \int_{-\infty}^{\infty} \tilde h = \sum_{k=1}^{\infty} 2^{-k}\int_{-\infty}^{\infty} g(x - r_k)\,dx = \sum_{k=1}^{\infty} 2^{-k}\,\|g\|_1 = \|g\|_1 = \frac{1}{\log 2} < \infty , \end{equation*}

using translation invariance again. Since \(\tilde h\) is nonnegative with finite integral, the set \(N = \{x : \tilde h(x) = \infty\}\) has measure \(0\) (if \(|N| > 0\) then \(\int \tilde h \ge c\,|N| \to \infty\) as \(c \to \infty\), a contradiction). Define

\begin{equation*} h(x) = \begin{cases} \tilde h(x) & \text{if } x \notin N,\\ 0 & \text{if } x \in N.\end{cases} \end{equation*}

Then \(h \colon \mathbb{R} \to \mathbb{R}\) is Borel measurable, \(h = \tilde h\) except on the set \(N\) of measure \(0\), and \(\int |h| = \int \tilde h < \infty\); thus \(h \in \mathcal{L}^1(\mathbb{R})\).

Conclusion. First we check that \(h^{*} = \tilde h^{*}\). Fix \(b \in \mathbb{R}\) and \(t > 0\) and put \(A = (b-t, b+t)\). Since \(h \ge 0\) we have \(|h| = h\), and \(h\) and \(\tilde h\) agree on \(A \setminus N\) while \(A \cap N\) has measure \(0\). The integral of a nonnegative measurable function over a set of measure \(0\) is \(0\): any simple function \(\sum_j c_j \chi_{E_j}\) lying below it has integral \(\sum_j c_j |E_j \cap A \cap N| = 0\) over that set, and the integral is the supremum of such numbers. Hence

\begin{equation*} \int_{A} \tilde h = \int_{A \setminus N} \tilde h + \int_{A \cap N} \tilde h = \int_{A\setminus N} h + 0 = \int_{A} h = \int_{b-t}^{b+t}|h| . \end{equation*}

As this holds for every \(b\) and every \(t > 0\), taking suprema over \(t\) gives \(h^{*} = \tilde h^{*}\). Fix \(b \in \mathbb{Q}\) and choose \(k\) with \(r_k = b\). Since \(\tilde h \ge 2^{-k} g(\cdot - r_k)\) pointwise, properties (i) and (ii) give

\begin{equation*} h^{*}(b) = \tilde h^{*}(b) \ \ge\ 2^{-k}\,\big(g(\cdot - r_k)\big)^{*}(b) = 2^{-k}\, g^{*}(b - r_k) = 2^{-k}\, g^{*}(0) = \infty . \end{equation*}

(Here we also used the obvious identity \((\lambda f)^{*} = \lambda f^{*}\) for \(\lambda > 0\).)

Hence \(h \in \mathcal{L}^1(\mathbb{R})\) and \(h^{*}(b) = \infty\) for every \(b \in \mathbb{Q}\), as required. Note that this is consistent with Exercise 12: the set \(\{b : h^{*}(b) = \infty\}\) contains \(\mathbb{Q}\) but still has measure \(0\).

Problem (4A.14)

Suppose \(h \in \mathcal{L}^1(\mathbb{R})\). Prove that

\begin{equation*} \big|\{b \in \mathbb{R} : h^{*}(b) \ge c\}\big| \le \frac{3}{c}\,\|h\|_1 \end{equation*}

for every \(c > 0\).

[This result slightly strengthens the Hardy–Littlewood maximal inequality (4.8) because the set on the left side above includes those \(b \in \mathbb{R}\) such that \(h^{*}(b) = c\). A much deeper strengthening comes from replacing the constant \(3\) in the Hardy–Littlewood maximal inequality with a smaller constant. In 2003, Antonios Melas answered what had been an open question about the best constant. He proved that the smallest constant that can replace \(3\) in the Hardy–Littlewood maximal inequality is \((11 + \sqrt{61})/12 \approx 1.56752\); see Annals of Mathematics 157 (2003), 647–688.]

Solution

Fix \(c > 0\) and set

\begin{equation*} E_c = \{b \in \mathbb{R} : h^{*}(b) \ge c\}, \qquad U_s = \{b \in \mathbb{R} : h^{*}(b) > s\} \ \ (s > 0). \end{equation*}

Each \(U_s\) is open by Exercise 9 in this section, and

\begin{equation*} E_c = \bigcap_{n > 1/c} U_{c - 1/n}, \end{equation*}

so \(E_c\) is a Borel set; in particular it is Lebesgue measurable. (As in Exercise 12, the argument below is equally valid with \(|\cdot|\) interpreted as outer measure, using monotonicity 2.5.)

Let \(n\) be any integer with \(n > 1/c\), so that \(c - \frac1n > 0\). If \(h^{*}(b) \ge c\) then \(h^{*}(b) > c - \frac1n\), so

\begin{equation*} E_c \subseteq U_{c - 1/n}. \end{equation*}

The Hardy–Littlewood maximal inequality (4.8), applied with the positive number \(c - \frac1n\) in place of \(c\), gives

\begin{equation*} \big|U_{c-1/n}\big| \ \le\ \frac{3}{c - \frac1n}\,\|h\|_1 . \end{equation*}

Hence, by monotonicity of Lebesgue measure,

\begin{equation*} |E_c| \ \le\ \frac{3\,\|h\|_1}{c - \frac1n} \qquad \text{for every integer } n > \tfrac1c . \end{equation*}

Since \(h \in \mathcal{L}^1(\mathbb{R})\) we have \(\|h\|_1 < \infty\), and \(\frac{3\|h\|_1}{c - 1/n} \to \frac{3\|h\|_1}{c}\) as \(n \to \infty\). Therefore

\begin{equation*} |E_c| \ \le\ \lim_{n \to \infty} \frac{3\,\|h\|_1}{c - \frac1n} \ = \ \frac{3}{c}\,\|h\|_1 , \end{equation*}

which is the desired inequality.

Exercises 4B

Problem (4B.1)

Suppose \(f \in L^1(\mathbb{R})\). Prove that

\begin{equation*} \lim_{t \downarrow 0} \frac{1}{2t} \int_{b-t}^{b+t} |f - f_{[b-t,\, b+t]}| = 0 \end{equation*}

for almost every \(b \in \mathbb{R}\).

[The book prefixes the exercises of this section with the following notation: for \(f \in L^1(\mathbb{R})\) and \(I\) an interval of \(\mathbb{R}\) with \(0 < |I| < \infty\), let \(f_I\) denote the average of \(f\) on \(I\). In other words, \(f_I = \frac{1}{|I|}\int_I f\).]

Solution

Let \(E\) be the set of \(b \in \mathbb{R}\) for which

\begin{equation*} \lim_{t \downarrow 0} \frac{1}{2t} \int_{b-t}^{b+t} |f - f(b)| = 0 . \end{equation*}

By the first version of the Lebesgue Differentiation Theorem (4.10), \(\mathbb{R} \setminus E\) has Lebesgue measure \(0\). We show that the desired conclusion holds at every \(b \in E\); this suffices.

So fix \(b \in E\) and fix \(t > 0\). Write \(I = [b-t,\, b+t]\), so that \(0 < |I| = 2t < \infty\) and hence \(f_I\) is defined, with

\begin{equation*} f_I = \frac{1}{2t}\int_{b-t}^{b+t} f . \end{equation*}

Note that \(f(b) \in \mathbb{R}\), because by definition (3.40 and 3.45) the elements of \(L^1(\mathbb{R})\) are real-valued functions. Note also that \(\int_{b-t}^{b+t}|f| \le \int_{-\infty}^{\infty}|f| < \infty\) and \(\int_{b-t}^{b+t}|f(b)| = 2t\,|f(b)| < \infty\), so every integral written below is the integral of a function that is integrable on \(I\); in particular all the quantities below are finite real numbers and the manipulations are legitimate.

First we estimate how far the average \(f_I\) is from the value \(f(b)\). Because \(\frac{1}{2t}\int_{b-t}^{b+t} f(b) = f(b)\), additivity of integration (3.21) gives

\begin{equation*} |f_I - f(b)| = \left| \frac{1}{2t}\int_{b-t}^{b+t} \bigl(f - f(b)\bigr) \right| \le \frac{1}{2t}\int_{b-t}^{b+t} |f - f(b)| , \end{equation*}

where the inequality holds because the absolute value of an integral is at most the integral of the absolute value (3.23).

Now the triangle inequality gives \(|f(x) - f_I| \le |f(x) - f(b)| + |f(b) - f_I|\) for every \(x \in I\). Integrating this over \(I\) (integration is order preserving, by 3.22, and additive, by 3.21) and using \(\int_I |f(b) - f_I| = 2t\,|f(b) - f_I|\), then dividing by \(2t\), we obtain

\begin{equation*} \frac{1}{2t}\int_{b-t}^{b+t} |f - f_I| \le \frac{1}{2t}\int_{b-t}^{b+t} |f - f(b)| + |f(b) - f_I| \le 2 \cdot \frac{1}{2t}\int_{b-t}^{b+t} |f - f(b)| , \end{equation*}

where the last inequality uses the estimate of the previous paragraph.

Because \(b \in E\), the right side of the last display has limit \(2 \cdot 0 = 0\) as \(t \downarrow 0\). Hence

\begin{equation*} \lim_{t \downarrow 0} \frac{1}{2t}\int_{b-t}^{b+t} |f - f_{[b-t,\, b+t]}| = 0 , \end{equation*}

as desired. Because this holds for every \(b \in E\) and \(\mathbb{R}\setminus E\) is a set of measure \(0\), the conclusion holds for almost every \(b \in \mathbb{R}\).

Problem (4B.2)

Suppose \(f \in L^1(\mathbb{R})\). Prove that

\begin{equation*} \lim_{t \downarrow 0} \sup\left\{ \frac{1}{|I|}\int_I |f - f_I| \; : \; I \text{ is an interval of length } t \text{ containing } b \right\} = 0 \end{equation*}

for almost every \(b \in \mathbb{R}\).

[As in the previous exercise, \(f_I\) denotes the average \(\frac{1}{|I|}\int_I f\) of \(f\) over an interval \(I\) with \(0 < |I| < \infty\).]

Solution

As in the previous exercise, let \(E\) denote the set of \(b \in \mathbb{R}\) such that

\begin{equation*} \lim_{t \downarrow 0} \frac{1}{2t}\int_{b-t}^{b+t} |f - f(b)| = 0 . \end{equation*}

The first version of the Lebesgue Differentiation Theorem (4.10) tells us that \(\mathbb{R}\setminus E\) has measure \(0\). We prove that the limit in the statement of this exercise equals \(0\) for every \(b \in E\).

Fix \(b \in E\). For each \(t > 0\) set

\begin{equation*} \Phi(t) = \sup\left\{ \frac{1}{|I|}\int_I |f - f_I| \; : \; I \text{ is an interval of length } t \text{ containing } b \right\} . \end{equation*}

The set over which this supremum is taken is nonempty (it contains the value coming from \(I = [b,\, b+t]\), for instance), and each \(I\) occurring in it has \(0 < |I| = t < \infty\), so \(f_I\) is defined for each such \(I\).

Now fix \(t > 0\) and let \(I\) be any interval of length \(t\) with \(b \in I\). Every \(x \in I\) satisfies \(|x - b| \le t\), because \(x\) and \(b\) both lie in \(I\) and \(I\) has length \(t\). Hence

\begin{equation*} I \subseteq [b-t,\, b+t] . \end{equation*}

Repeating the computation from the previous exercise with the interval \(I\) in place of \([b-t,\,b+t]\), we get

\begin{equation*} |f_I - f(b)| = \left|\frac{1}{|I|}\int_I \bigl(f - f(b)\bigr)\right| \le \frac{1}{|I|}\int_I |f - f(b)| \end{equation*}

by 3.23, and therefore, by the triangle inequality applied pointwise and then integrated over \(I\),

\begin{equation*} \frac{1}{|I|}\int_I |f - f_I| \le \frac{1}{|I|}\int_I |f - f(b)| + |f(b) - f_I| \le \frac{2}{|I|}\int_I |f - f(b)| . \end{equation*}

Because \(|f - f(b)| \ge 0\) and \(I \subseteq [b-t,\,b+t]\), we have \(\int_I |f - f(b)| \le \int_{b-t}^{b+t} |f - f(b)|\). Since \(|I| = t\), combining this with the previous display yields

\begin{equation*} \frac{1}{|I|}\int_I |f - f_I| \le \frac{2}{t}\int_{b-t}^{b+t} |f - f(b)| = 4 \cdot \frac{1}{2t}\int_{b-t}^{b+t}|f - f(b)| . \end{equation*}

The right side of the last display does not depend on \(I\); it depends only on \(b\) and \(t\). Thus it is an upper bound for the set whose supremum defines \(\Phi(t)\), and hence

\begin{equation*} 0 \le \Phi(t) \le 4 \cdot \frac{1}{2t}\int_{b-t}^{b+t}|f - f(b)| . \end{equation*}

Because \(b \in E\), the right side has limit \(0\) as \(t \downarrow 0\). By the squeeze principle, \(\lim_{t \downarrow 0} \Phi(t) = 0\), which is exactly the assertion of the exercise at the point \(b\).

Because \(\mathbb{R}\setminus E\) has measure \(0\), this holds for almost every \(b \in \mathbb{R}\).

Problem (4B.3)

Suppose \(f : \mathbb{R} \to \mathbb{R}\) is a Lebesgue measurable function such that \(f^2 \in L^1(\mathbb{R})\). Prove that

\begin{equation*} \lim_{t \downarrow 0} \frac{1}{2t}\int_{b-t}^{b+t} |f - f(b)|^2 = 0 \end{equation*}

for almost every \(b \in \mathbb{R}\).

Solution

The hypothesis here is that \(f^2 \in L^1(\mathbb{R})\), not that \(f \in L^1(\mathbb{R})\), so 4.10 does not apply directly to \(f\) and we must build the conclusion out of functions to which it does apply. We first record a localized form of 4.10, which frees us from having to integrate over all of \(\mathbb{R}\).

Localization Lemma. Suppose \(h : \mathbb{R} \to \mathbb{R}\) is Lebesgue measurable and \(\int_{-n}^{n} |h| < \infty\) for every \(n \in \mathbb{Z}^+\). Then

\begin{equation*} \lim_{t \downarrow 0} \frac{1}{2t}\int_{b-t}^{b+t} |h - h(b)| = 0 \end{equation*}

for almost every \(b \in \mathbb{R}\); in particular \(\frac{1}{2t}\int_{b-t}^{b+t} h \to h(b)\) as \(t \downarrow 0\) for almost every \(b \in \mathbb{R}\).

Proof of the lemma. For \(n \in \mathbb{Z}^+\) let \(h_n = h \chi_{(-n,\, n)}\). Then \(h_n\) is Lebesgue measurable and \(\int |h_n| = \int_{-n}^{n}|h| < \infty\), so \(h_n \in L^1(\mathbb{R})\). By 4.10 there is a set \(N_n \subseteq \mathbb{R}\) with \(|N_n| = 0\) such that

\begin{equation*} \lim_{t\downarrow 0} \frac{1}{2t}\int_{b-t}^{b+t}|h_n - h_n(b)| = 0 \quad \text{for all } b \in \mathbb{R}\setminus N_n . \end{equation*}

Let \(N = \bigcup_{n=1}^{\infty} N_n\); then \(|N| = 0\) by countable subadditivity (2.8). Suppose \(b \in \mathbb{R}\setminus N\). Choose \(n \in \mathbb{Z}^+\) with \(|b| < n\) and put \(\delta = n - |b| > 0\). If \(0 < t < \delta\) then \((b-t,\, b+t) \subseteq (-n,\, n)\), so \(h = h_n\) on \((b-t,\,b+t)\) and \(h(b) = h_n(b)\); hence

\begin{equation*} \frac{1}{2t}\int_{b-t}^{b+t}|h - h(b)| = \frac{1}{2t}\int_{b-t}^{b+t}|h_n - h_n(b)| . \end{equation*}

Because \(b \notin N_n\), the right side has limit \(0\) as \(t \downarrow 0\), and therefore so does the left side. The last assertion of the lemma follows because

\begin{equation*} \left|\frac{1}{2t}\int_{b-t}^{b+t} h - h(b)\right| = \left|\frac{1}{2t}\int_{b-t}^{b+t}\bigl(h - h(b)\bigr)\right| \le \frac{1}{2t}\int_{b-t}^{b+t}|h - h(b)| , \end{equation*}

using 3.23. This completes the proof of the lemma.

Now we turn to the exercise. For \(q \in \mathbb{Q}\) define \(h_q : \mathbb{R}\to\mathbb{R}\) by

\begin{equation*} h_q(x) = \bigl(f(x) - q\bigr)^2 . \end{equation*}

Each \(h_q\) is Lebesgue measurable, because sums and products of measurable functions are measurable (2.46). Furthermore, for each \(n \in \mathbb{Z}^+\) the inequality \((\alpha - \beta)^2 \le 2\alpha^2 + 2\beta^2\) gives

\begin{equation*} \int_{-n}^{n} h_q \le 2\int_{-n}^{n} f^2 + 2q^2 \cdot 2n \le 2\int_{-\infty}^{\infty} f^2 + 4nq^2 < \infty , \end{equation*}

because \(f^2 \in L^1(\mathbb{R})\). Since \(h_q \ge 0\), we have \(|h_q| = h_q\), so the display above says \(\int_{-n}^{n}|h_q| < \infty\) for every \(n \in \mathbb{Z}^+\); thus \(h_q\) satisfies the hypothesis of the Localization Lemma. Let \(N_q\) be a set of measure \(0\) such that

\begin{equation*} \lim_{t\downarrow 0}\frac{1}{2t}\int_{b-t}^{b+t} h_q = h_q(b) = \bigl(f(b) - q\bigr)^2 \quad \text{for all } b \in \mathbb{R}\setminus N_q , \end{equation*}

which the lemma provides. Let

\begin{equation*} N = \bigcup_{q \in \mathbb{Q}} N_q . \end{equation*}

Because \(\mathbb{Q}\) is countable, \(|N| = 0\).

Suppose now \(b \in \mathbb{R}\setminus N\); we show the conclusion of the exercise holds at \(b\). Let \(\varepsilon > 0\). Because \(f(b) \in \mathbb{R}\) and \(\mathbb{Q}\) is dense in \(\mathbb{R}\), there exists \(q \in \mathbb{Q}\) with \(|f(b) - q| < \varepsilon\). For every \(x \in \mathbb{R}\),

\begin{equation*} |f(x) - f(b)|^2 \le \bigl(|f(x) - q| + |q - f(b)|\bigr)^2 \le 2\bigl(f(x)-q\bigr)^2 + 2\bigl(q - f(b)\bigr)^2 , \end{equation*}

using \((\alpha+\beta)^2 \le 2\alpha^2 + 2\beta^2\). Integrating over \((b-t,\,b+t)\) and dividing by \(2t\) gives

\begin{equation*} \frac{1}{2t}\int_{b-t}^{b+t}|f - f(b)|^2 \le 2 \cdot \frac{1}{2t}\int_{b-t}^{b+t} h_q + 2\bigl(q - f(b)\bigr)^2 . \end{equation*}

Letting \(t \downarrow 0\) and using \(b \notin N_q\), we obtain

\begin{equation*} \limsup_{t\downarrow 0}\; \frac{1}{2t}\int_{b-t}^{b+t}|f - f(b)|^2 \le 2\bigl(f(b)-q\bigr)^2 + 2\bigl(q-f(b)\bigr)^2 = 4\bigl(f(b)-q\bigr)^2 < 4\varepsilon^2 . \end{equation*}

The number \(\varepsilon > 0\) was arbitrary and the left side of the last display does not depend on \(\varepsilon\); hence that limit superior is \(0\). Since \(\frac{1}{2t}\int_{b-t}^{b+t}|f-f(b)|^2 \ge 0\) for every \(t > 0\), the corresponding limit inferior is also \(0\), and therefore

\begin{equation*} \lim_{t\downarrow 0}\frac{1}{2t}\int_{b-t}^{b+t}|f - f(b)|^2 = 0 . \end{equation*}

As \(|N| = 0\), this holds for almost every \(b \in \mathbb{R}\).

Problem (4B.4)

Prove that the Lebesgue Differentiation Theorem (4.19) still holds if the hypothesis that \(\int_{-\infty}^{\infty} |f| < \infty\) is weakened to the requirement that \(\int_{-\infty}^{x} |f| < \infty\) for all \(x \in \mathbb{R}\).

Solution

Suppose \(f : \mathbb{R}\to\mathbb{R}\) is Lebesgue measurable and

\begin{equation*} \int_{-\infty}^{x} |f| < \infty \quad \text{for every } x \in \mathbb{R} . \end{equation*}

Then \(f\chi_{(-\infty,\, x)} \in L^1(\mathbb{R})\) for every \(x\), so the function \(g : \mathbb{R}\to\mathbb{R}\) defined by

\begin{equation*} g(x) = \int_{-\infty}^{x} f \end{equation*}

is well defined and real valued. We must show that \(g\) is differentiable at \(b\) with \(g^{\prime}(b) = f(b)\) for almost every \(b \in \mathbb{R}\).

Fix \(n \in \mathbb{Z}^+\) and define \(f_n : \mathbb{R}\to\mathbb{R}\) by

\begin{equation*} f_n = f \chi_{(-\infty,\, n)} . \end{equation*}

Then \(f_n\) is Lebesgue measurable and

\begin{equation*} \int_{-\infty}^{\infty} |f_n| = \int_{-\infty}^{n} |f| < \infty , \end{equation*}

so \(f_n \in L^1(\mathbb{R})\). Define \(g_n : \mathbb{R}\to\mathbb{R}\) by \(g_n(x) = \int_{-\infty}^{x} f_n\). By the Lebesgue Differentiation Theorem in the form 4.19, applied to \(f_n\), there is a set \(N_n \subseteq \mathbb{R}\) with \(|N_n| = 0\) such that

\begin{equation*} g_n^{\prime}(b) = f_n(b) \quad \text{for every } b \in \mathbb{R}\setminus N_n . \end{equation*}

Now observe that if \(x < n\) then \(f = f_n\) on \((-\infty,\, x)\), and hence

\begin{equation*} g_n(x) = \int_{-\infty}^{x} f_n = \int_{-\infty}^{x} f = g(x) . \end{equation*}

Thus \(g\) and \(g_n\) agree on the open interval \((-\infty,\, n)\). Suppose \(b < n\). If \(|t|\) is small enough that \(b + t < n\) and \(t \ne 0\), then

\begin{equation*} \frac{g(b+t) - g(b)}{t} = \frac{g_n(b+t) - g_n(b)}{t} , \end{equation*}

because both \(b\) and \(b+t\) lie in \((-\infty,\, n)\). Hence the two difference quotients in the definition of the derivative (4.16) agree for all \(t\) near \(0\), so one of the limits exists if and only if the other does, and then they are equal. In other words: if \(b < n\) and \(g_n\) is differentiable at \(b\), then \(g\) is differentiable at \(b\) and \(g^{\prime}(b) = g_n^{\prime}(b)\). Moreover \(f_n(b) = f(b)\) for \(b < n\), directly from the definition of \(f_n\).

Combining the two previous paragraphs: if \(b \in (-\infty,\, n)\setminus N_n\), then \(g\) is differentiable at \(b\) and

\begin{equation*} g^{\prime}(b) = g_n^{\prime}(b) = f_n(b) = f(b) . \end{equation*}

Let \(N = \bigcup_{n=1}^{\infty} N_n\). By countable subadditivity of outer measure (2.8), \(|N| = 0\). Suppose \(b \in \mathbb{R}\setminus N\). Choose \(n \in \mathbb{Z}^+\) with \(n > b\). Because \(b \notin N_n\) and \(b < n\), the previous paragraph gives \(g^{\prime}(b) = f(b)\).

Hence \(g^{\prime}(b) = f(b)\) for every \(b \in \mathbb{R}\setminus N\), and thus for almost every \(b \in \mathbb{R}\). This proves that 4.19 holds under the weakened hypothesis.

Problem (4B.5)

Suppose \(f : \mathbb{R}\to\mathbb{R}\) is a Lebesgue measurable function. Prove that

\begin{equation*} |f(b)| \le f^*(b) \end{equation*}

for almost every \(b \in \mathbb{R}\).

Solution

Recall from 4.6 that the Hardy–Littlewood maximal function of \(f\) is

\begin{equation*} f^*(b) = \sup_{t > 0} \frac{1}{2t}\int_{b-t}^{b+t}|f| , \end{equation*}

with values in \([0,\infty]\). Note that no integrability hypothesis is assumed on \(f\), so we must handle the points near which \(|f|\) fails to be integrable; at those points the inequality will turn out to hold trivially because \(f^*\) is infinite there.

Let

\begin{equation*} \mathcal{F} = \left\{ (p,q) \; : \; p, q \in \mathbb{Q},\; p < q, \; \int_p^q |f| < \infty \right\} , \end{equation*}

which is a subset of \(\mathbb{Q}\times\mathbb{Q}\) and hence countable. Let

\begin{equation*} L = \bigcup_{(p,q) \in \mathcal{F}} (p,q) . \end{equation*}

So \(L\) is an open subset of \(\mathbb{R}\). We treat \(\mathbb{R}\setminus L\) and \(L\) separately.

Step 1: the inequality holds at every \(b \in \mathbb{R}\setminus L\). Suppose \(b \notin L\) and let \(t > 0\). If we had \(\int_{b-t}^{b+t}|f| < \infty\), then choosing rational numbers \(p, q\) with \(b - t < p < b < q < b + t\) would give \(\int_p^q |f| \le \int_{b-t}^{b+t}|f| < \infty\) (the inequality holds because \(|f| \ge 0\) and \((p,q) \subseteq (b-t,\,b+t)\)), so \((p,q) \in \mathcal{F}\) and therefore \(b \in (p,q) \subseteq L\), a contradiction. Hence

\begin{equation*} \int_{b-t}^{b+t}|f| = \infty \quad \text{for every } t > 0 , \end{equation*}

and thus \(f^*(b) = \infty\). Because \(f\) is real valued, \(|f(b)| < \infty = f^*(b)\). So the desired inequality holds at every point of \(\mathbb{R}\setminus L\), with no exceptional set at all.

Step 2: the inequality holds at almost every \(b \in L\). Fix \((p,q) \in \mathcal{F}\) and define

\begin{equation*} h = |f| \chi_{(p,\, q)} . \end{equation*}

Then \(h\) is Lebesgue measurable, \(h \ge 0\), and \(\int_{-\infty}^{\infty} h = \int_p^q |f| < \infty\), so \(h \in L^1(\mathbb{R})\). By the first version of the Lebesgue Differentiation Theorem (4.10) there is a set \(N_{p,q}\) with \(|N_{p,q}| = 0\) such that

\begin{equation*} \lim_{t\downarrow 0}\frac{1}{2t}\int_{b-t}^{b+t}|h - h(b)| = 0 \quad \text{for all } b \in \mathbb{R}\setminus N_{p,q} . \end{equation*}

Because

\begin{equation*} \left|\frac{1}{2t}\int_{b-t}^{b+t} h \;-\; h(b)\right| \le \frac{1}{2t}\int_{b-t}^{b+t}|h - h(b)| \end{equation*}

by 3.23, we get

\begin{equation*} \lim_{t\downarrow 0}\frac{1}{2t}\int_{b-t}^{b+t} h = h(b) \quad \text{for all } b \in \mathbb{R}\setminus N_{p,q} . \end{equation*}

Now suppose \(b \in (p,q)\setminus N_{p,q}\), and let \(\delta = \min\{b - p,\; q - b\} > 0\). If \(0 < t < \delta\) then \((b-t,\,b+t) \subseteq (p,q)\), so \(h = |f|\) on \((b-t,\,b+t)\); also \(h(b) = |f(b)|\). Consequently

\begin{equation*} \frac{1}{2t}\int_{b-t}^{b+t}|f| = \frac{1}{2t}\int_{b-t}^{b+t} h \longrightarrow h(b) = |f(b)| \quad \text{as } t \downarrow 0 . \end{equation*}

By the definition of \(f^*\), every \(t > 0\) satisfies \(\frac{1}{2t}\int_{b-t}^{b+t}|f| \le f^*(b)\). A limit of quantities each of which is at most \(f^*(b)\) is itself at most \(f^*(b)\), so the display above yields

\begin{equation*} |f(b)| = \lim_{t\downarrow 0}\frac{1}{2t}\int_{b-t}^{b+t}|f| \le f^*(b) . \end{equation*}

Finally, let

\begin{equation*} N = \bigcup_{(p,q)\in\mathcal{F}} N_{p,q} . \end{equation*}

Because \(\mathcal{F}\) is countable, \(|N| = 0\) by countable subadditivity of outer measure (2.8). If \(b \in \mathbb{R}\setminus N\), then either \(b \notin L\), in which case Step 1 gives \(|f(b)| \le f^*(b)\), or \(b \in L\), in which case \(b \in (p,q)\) for some \((p,q) \in \mathcal{F}\) and \(b \notin N_{p,q}\), so Step 2 gives \(|f(b)| \le f^*(b)\).

Hence \(|f(b)| \le f^*(b)\) for every \(b \in \mathbb{R}\setminus N\), and therefore for almost every \(b \in \mathbb{R}\).

Problem (4B.6)

Prove that if \(h \in L^1(\mathbb{R})\) and \(\int_{-\infty}^{s} h = 0\) for all \(s \in \mathbb{R}\), then \(h(s) = 0\) for almost every \(s \in \mathbb{R}\).

Solution

Suppose \(h \in L^1(\mathbb{R})\) and \(\int_{-\infty}^{s} h = 0\) for every \(s \in \mathbb{R}\). Define \(g : \mathbb{R}\to\mathbb{R}\) by

\begin{equation*} g(x) = \int_{-\infty}^{x} h , \end{equation*}

which makes sense because \(h \in L^1(\mathbb{R})\). The hypothesis says precisely that

\begin{equation*} g(x) = 0 \quad \text{for every } x \in \mathbb{R} , \end{equation*}

so \(g\) is the zero function. In particular \(g\) is differentiable at every \(b \in \mathbb{R}\) with

\begin{equation*} g^{\prime}(b) = \lim_{t\to 0}\frac{g(b+t) - g(b)}{t} = \lim_{t\to 0}\frac{0 - 0}{t} = 0 . \end{equation*}

On the other hand, the Lebesgue Differentiation Theorem in its second version (4.19), applied to \(h \in L^1(\mathbb{R})\) and to this same function \(g\), states that

\begin{equation*} g^{\prime}(b) = h(b) \quad \text{for almost every } b \in \mathbb{R} . \end{equation*}

Comparing the two previous displays, we get \(h(b) = g^{\prime}(b) = 0\) for almost every \(b \in \mathbb{R}\). That is, \(h(s) = 0\) for almost every \(s \in \mathbb{R}\), as desired.

Problem (4B.7)

Give an example of a Borel subset of \(\mathbb{R}\) whose density at \(0\) is not defined.

Solution

Recall from 4.22 that the density of a set \(E \subseteq \mathbb{R}\) at \(0\) is

\begin{equation*} \lim_{t\downarrow 0}\frac{|E \cap (-t,\, t)|}{2t} \end{equation*}

if this limit exists, and is undefined otherwise. We produce a Borel set \(E\) for which the ratio above oscillates between \(\frac{1}{5}\) and \(\frac{4}{5}\) as \(t \downarrow 0\).

Define

\begin{equation*} E = \bigcup_{k=0}^{\infty} \Bigl( \bigl[-4^{-2k},\, -4^{-2k-1}\bigr] \cup \bigl[4^{-2k-1},\, 4^{-2k}\bigr] \Bigr) . \end{equation*}

Thus \(E\) is symmetric about \(0\), and \(E \cap (0,\infty)\) consists of the closed intervals

\begin{equation*} [4^{-1},\, 4^{0}],\quad [4^{-3},\, 4^{-2}],\quad [4^{-5},\, 4^{-4}],\quad \dots , \end{equation*}

which pile up at \(0\) with gaps between them. Being a countable union of closed intervals, \(E\) is a Borel set.

These intervals are pairwise disjoint: for \(j \ge 0\) the interval indexed by \(j+1\) has right endpoint \(4^{-2j-2}\), which is strictly less than the left endpoint \(4^{-2j-1}\) of the interval indexed by \(j\). So countable additivity of Lebesgue measure applies to any subfamily of them.

We first record the basic length computation. For \(j \ge 0\),

\begin{equation*} \bigl| [4^{-2j-1},\, 4^{-2j}] \bigr| = 4^{-2j} - 4^{-2j-1} = \tfrac{3}{4}\,4^{-2j} = \tfrac{3}{4}\,16^{-j} . \end{equation*}

Case 1: \(t = t_m := 4^{-2m}\) for \(m \in \mathbb{Z}^+\). Then

\begin{equation*} E \cap (0,\, t_m) = \bigl[4^{-2m-1},\, 4^{-2m}\bigr) \cup \bigcup_{j > m} \bigl[4^{-2j-1},\, 4^{-2j}\bigr] , \end{equation*}

because the intervals of \(E\) lying in \((0,\infty)\) that meet \((0, 4^{-2m})\) are exactly \([4^{-2j-1}, 4^{-2j}]\) for \(j \ge m\) (those with \(j < m\) lie entirely in \([4^{-2m+1},\infty)\), to the right of \(t_m\)), and the one with \(j = m\) has its right endpoint at \(t_m\), which is excluded. Since a single point has measure \(0\), countable additivity gives

\begin{equation*} |E \cap (0,\, t_m)| = \sum_{j \ge m} \tfrac{3}{4}\, 16^{-j} = \tfrac{3}{4}\cdot \frac{16^{-m}}{1 - \frac{1}{16}} = \tfrac{3}{4}\cdot\frac{16}{15}\cdot 16^{-m} = \tfrac{4}{5}\, 4^{-2m} = \tfrac{4}{5}\, t_m . \end{equation*}

By the symmetry of \(E\) about \(0\) (and since \(\{0\}\) has measure \(0\)),

\begin{equation*} \frac{|E \cap (-t_m,\, t_m)|}{2t_m} = \frac{2 \cdot \frac{4}{5}t_m}{2t_m} = \frac{4}{5} . \end{equation*}

Case 2: \(t = s_m := 4^{-2m-1}\) for \(m \in \mathbb{Z}^+\). The point \(s_m\) is the left endpoint of the interval \([4^{-2m-1}, 4^{-2m}]\) of \(E\), so that interval meets \((0, s_m)\) in the empty set, and the intervals with \(j < m\) lie even further to the right. The intervals with \(j \ge m+1\) have right endpoint \(4^{-2j} \le 4^{-2m-2} < s_m\), so they lie entirely inside \((0, s_m)\). Hence \(E \cap (0, s_m)\) is exactly the disjoint union of \([4^{-2j-1}, 4^{-2j}]\) over \(j \ge m+1\), and

\begin{equation*} |E \cap (0,\, s_m)| = \sum_{j \ge m+1} \tfrac{3}{4}\,16^{-j} = \tfrac{4}{5}\,16^{-(m+1)} = \tfrac{4}{5}\cdot \frac{4^{-2m}}{16} = \tfrac{1}{5}\cdot \frac{4^{-2m}}{4} = \tfrac{1}{5}\, s_m , \end{equation*}

where the second equality uses the computation of Case 1 with \(m\) replaced by \(m+1\). Again by symmetry,

\begin{equation*} \frac{|E \cap (-s_m,\, s_m)|}{2s_m} = \frac{2\cdot\frac{1}{5}s_m}{2s_m} = \frac{1}{5} . \end{equation*}

Both \(t_m \downarrow 0\) and \(s_m \downarrow 0\) as \(m \to \infty\), yet along the first sequence the ratio \(\frac{|E\cap(-t,t)|}{2t}\) is constantly \(\frac{4}{5}\) and along the second it is constantly \(\frac{1}{5}\). Because \(\frac{4}{5} \ne \frac{1}{5}\), the limit

\begin{equation*} \lim_{t\downarrow 0}\frac{|E \cap (-t,\, t)|}{2t} \end{equation*}

does not exist. Therefore the density of the Borel set \(E\) at \(0\) is not defined.

This is consistent with the Lebesgue Density Theorem (4.24): that theorem asserts that a Lebesgue measurable set has density \(1\) at almost every one of its points and density \(0\) at almost every point of its complement, and here the single point \(0\) (a set of measure \(0\)) is exceptional.

Problem (4B.8)

Give an example of a Borel subset of \(\mathbb{R}\) whose density at \(0\) is \(\frac{1}{3}\).

Solution

Recall (4.22) that the density of a set \(A \subseteq \mathbb{R}\) at \(0\) is

\begin{equation*} \lim_{t \downarrow 0} \frac{|A \cap (-t, t)|}{2t}, \end{equation*}

if this limit exists.

Two remarks explain the shape of the example below; neither is needed for the proof. First, one cannot simply demand that the set occupy exactly the proportion \(\frac{1}{3}\) of every interval \((0,t)\): by 4.20 there is no Lebesgue measurable \(A \subseteq [0,1]\) with \(|A \cap [0,b]| = \frac{b}{2}\) for all \(b \in [0,1]\), and the proof of 4.20 works verbatim with \(\frac{b}{2}\) replaced by \(cb\) for any \(c \in (0,1)\), since the derivative it produces is then the constant \(c\), a value never taken by a characteristic function. Second, a set built on a geometric scale, occupying a fixed proportion of each interval \(\big[2^{-k-1}, 2^{-k}\big)\), has a ratio \(\frac{|A \cap (0,t)|}{t}\) that is unchanged under \(t \mapsto t/2\) and is not constant in \(t\), hence has no limit as \(t \downarrow 0\); that is the mechanism behind Exercise 7. The construction below instead puts blocks on the harmonic scale \(\frac{1}{k}\), with the \(k\)-th block of length \(\frac{1}{3k(k+1)}\), small compared with its distance \(\frac{1}{k+1}\) from the origin. Truncating in the middle of a block then perturbs the ratio by an amount tending to \(0\), and the limit exists.

Construction. For \(k \in \mathbb{Z}^{+}\) let

\begin{equation*} J_k = \Big[\tfrac{1}{k+1}, \; \tfrac{1}{k+1} + \tfrac{1}{3k(k+1)}\Big), \end{equation*}

which is the left third of the interval \(\big[\frac{1}{k+1}, \frac{1}{k}\big)\), because

\begin{equation*} \frac{1}{k} - \frac{1}{k+1} = \frac{1}{k(k+1)}. \end{equation*}

In particular the intervals \(J_1, J_2, \ldots\) are disjoint, \(J_k \subseteq \big[\frac{1}{k+1}, \frac{1}{k}\big) \subseteq (0,1)\), and \(|J_k| = \frac{1}{3k(k+1)}\). Put

\begin{equation*} F = \bigcup_{k=1}^{\infty} J_k \quad \text{and} \quad E = F \cup (-F), \end{equation*}

where \(-F = \{-x : x \in F\}\). Each \(J_k\) is an interval, hence a Borel set, so \(F\) is a Borel set (a countable union of Borel sets), and \(-F\) is Borel because \(x \mapsto -x\) is a homeomorphism of \(\mathbb{R}\). Thus \(E\) is a Borel subset of \(\mathbb{R}\).

The density of \(E\) at \(0\) equals \(\frac{1}{3}\). Because \(F \subseteq (0,1)\) and \(-F \subseteq (-1,0)\), the point \(0\) lies in neither, so for every \(t > 0\) the set \(E \cap (-t,t)\) is the union of the two disjoint Borel sets \(F \cap (0,t)\) and \((-F) \cap (-t,0)\). Since outer measure restricted to the Borel sets is a measure (2.68), it is additive on that pair, so

\begin{equation*} |E \cap (-t,t)| = |F \cap (0,t)| + |(-F) \cap (-t, 0)| = 2\,|F \cap (0,t)|, \end{equation*}

the second equality because Lebesgue measure is invariant under \(x \mapsto -x\): reflection maps an open interval onto an open interval of the same length, so it carries any cover of a set by open intervals to a cover of the reflected set by open intervals of the same total length, and hence preserves outer measure (this is the argument used for translation invariance in 2.7). Hence

\begin{equation*} \frac{|E \cap (-t,t)|}{2t} = \frac{|F \cap (0,t)|}{t}, \end{equation*}

so it suffices to show that \(\frac{|F \cap (0,t)|}{t} \to \frac{1}{3}\) as \(t \downarrow 0\).

First compute \(|F \cap (0,t)|\) at the special values \(t = \frac{1}{n}\), where \(n \in \mathbb{Z}^{+}\). If \(k \geq n\) then \(J_k \subseteq \big[\frac{1}{k+1}, \frac{1}{k}\big) \subseteq \big(0, \frac{1}{n}\big)\), while if \(k < n\) then every element of \(J_k\) is at least \(\frac{1}{k+1} \geq \frac{1}{n}\), so \(J_k\) is disjoint from \(\big(0,\frac{1}{n}\big)\). Hence \(F \cap \big(0,\frac{1}{n}\big) = \bigcup_{k \geq n} J_k\), a disjoint union of Borel sets, and countable additivity (2.68) together with the formula for the outer measure of an interval (2.14) gives

\begin{equation*} \Big|F \cap \big(0, \tfrac{1}{n}\big)\Big| = \sum_{k=n}^{\infty} |J_k| = \frac{1}{3}\sum_{k=n}^{\infty}\Big(\frac{1}{k} - \frac{1}{k+1}\Big) = \frac{1}{3n}, \end{equation*}

the last equality by telescoping.

Now let \(t \in (0,1]\) and put \(n = \lfloor 1/t \rfloor\), so that \(n \in \mathbb{Z}^{+}\) and \(\frac{1}{n+1} < t \leq \frac{1}{n}\). Since \(\big(0,\frac{1}{n+1}\big) \subseteq (0,t) \subseteq \big(0,\frac{1}{n}\big]\) and the single point \(\frac{1}{n}\) has outer measure \(0\), monotonicity of outer measure gives

\begin{equation*} \frac{1}{3(n+1)} = \Big|F \cap \big(0,\tfrac{1}{n+1}\big)\Big| \leq |F \cap (0,t)| \leq \Big|F \cap \big(0,\tfrac{1}{n}\big)\Big| = \frac{1}{3n}. \end{equation*}

Combining this with \(\frac{1}{n+1} < t \leq \frac{1}{n}\) we obtain

\begin{equation*} \frac{n}{3(n+1)} = \frac{1/\big(3(n+1)\big)}{1/n} \leq \frac{|F \cap (0,t)|}{t} \leq \frac{1/(3n)}{1/(n+1)} = \frac{n+1}{3n}. \end{equation*}

Both outer bounds are functions of \(n\) alone, the left one increasing in \(n\) and the right one decreasing in \(n\), and both tend to \(\frac{1}{3}\) as \(n \to \infty\). Given \(\varepsilon > 0\), pick \(N \in \mathbb{Z}^{+}\) with \(\frac{N}{3(N+1)} > \frac{1}{3} - \varepsilon\) and \(\frac{N+1}{3N} < \frac{1}{3} + \varepsilon\). If \(0 < t \leq \frac{1}{N}\) then \(n = \lfloor 1/t \rfloor \geq N\), so the displayed bounds give \(\big|\frac{|F \cap (0,t)|}{t} - \frac{1}{3}\big| < \varepsilon\). Therefore

\begin{equation*} \lim_{t \downarrow 0} \frac{|F \cap (0,t)|}{t} = \frac{1}{3}, \end{equation*}

and consequently the density of the Borel set \(E\) at \(0\) equals \(\frac{1}{3}\).

Problem (4B.9)

Prove that if \(t \in [0,1]\), then there exists a Borel set \(E \subseteq \mathbb{R}\) such that the density of \(E\) at \(0\) is \(t\).

Solution

To avoid a clash with the variable \(t\) appearing in the definition of density (4.22), write the given number as \(c\); thus \(c \in [0,1]\) is given and we must produce a Borel set \(E \subseteq \mathbb{R}\) whose density at \(0\) equals \(c\), meaning

\begin{equation*} \lim_{s \downarrow 0} \frac{|E \cap (-s, s)|}{2s} = c. \end{equation*}

The construction is the one used in Exercise 8, with the proportion \(\frac{1}{3}\) replaced by \(c\). For \(k \in \mathbb{Z}^{+}\) let

\begin{equation*} J_k = \Big[\tfrac{1}{k+1}, \; \tfrac{1}{k+1} + \tfrac{c}{k(k+1)}\Big). \end{equation*}

Because \(\frac{1}{k} - \frac{1}{k+1} = \frac{1}{k(k+1)}\) and \(0 \leq c \leq 1\), the interval \(J_k\) is the left-hand fraction \(c\) of the interval \(\big[\frac{1}{k+1}, \frac{1}{k}\big)\); in particular

\begin{equation*} J_k \subseteq \Big[\tfrac{1}{k+1}, \tfrac{1}{k}\Big) \subseteq (0,1), \qquad |J_k| = \frac{c}{k(k+1)}, \end{equation*}

and the intervals \(J_1, J_2, \ldots\) are disjoint because the intervals \(\big[\frac{1}{k+1}, \frac{1}{k}\big)\) are disjoint. (If \(c = 0\) then each \(J_k\) is empty; if \(c = 1\) then \(J_k = \big[\frac{1}{k+1}, \frac{1}{k}\big)\). Both extreme cases are allowed in what follows.) Set

\begin{equation*} F = \bigcup_{k=1}^{\infty} J_k, \qquad E = F \cup (-F), \end{equation*}

where \(-F = \{-x : x \in F\}\). Each \(J_k\) is an interval and hence Borel, so \(F\) is Borel as a countable union of Borel sets; \(-F\) is Borel because \(x \mapsto -x\) is a homeomorphism of \(\mathbb{R}\) and so maps Borel sets to Borel sets. Thus \(E\) is a Borel subset of \(\mathbb{R}\).

Step 1: reduction to a one-sided statement. Since \(F \subseteq (0,1)\) and \(-F \subseteq (-1,0)\), the point \(0\) belongs to neither, so for each \(s > 0\) the sets \(F \cap (0,s)\) and \((-F) \cap (-s,0)\) are disjoint Borel sets whose union is \(E \cap (-s,s)\). Outer measure restricted to the Borel sets is a measure (2.68), hence additive on that pair. Moreover Lebesgue measure is invariant under the reflection \(x \mapsto -x\): reflection maps an open interval onto an open interval of the same length, so it carries any cover of a set by open intervals to a cover of the reflected set by open intervals of the same total length, and therefore preserves outer measure (this is the argument used for translation invariance in 2.7). Since \((-F) \cap (-s,0)\) is the reflection of \(F \cap (0,s)\), we get

\begin{equation*} |E \cap (-s,s)| = |F \cap (0,s)| + |F \cap (0,s)| = 2\,|F \cap (0,s)|, \end{equation*}

so that

\begin{equation*} \frac{|E \cap (-s,s)|}{2s} = \frac{|F \cap (0,s)|}{s}. \end{equation*}

It therefore suffices to prove

\begin{equation*} \lim_{s \downarrow 0} \frac{|F \cap (0,s)|}{s} = c. \end{equation*}

Step 2: the measure of \(F\) at the scales \(\frac{1}{n}\). Fix \(n \in \mathbb{Z}^{+}\). Since \(J_k \subseteq \big(0, \frac{1}{n}\big)\) for \(k \geq n\), while \(J_k \subseteq \big[\frac{1}{k+1}, \frac{1}{k}\big)\) is disjoint from \(\big(0, \frac{1}{n}\big)\) for \(k < n\), we have \(F \cap \big(0,\frac{1}{n}\big) = \bigcup_{k \geq n} J_k\). The \(J_k\) are disjoint Borel sets, so countable additivity of outer measure on Borel sets (2.68), together with the fact that the outer measure of an interval is its length (2.14), gives

\begin{equation*} \Big|F \cap \big(0,\tfrac{1}{n}\big)\Big| = \sum_{k=n}^{\infty} \frac{c}{k(k+1)} = c \sum_{k=n}^{\infty} \Big(\frac{1}{k} - \frac{1}{k+1}\Big) = \frac{c}{n}, \end{equation*}

the last equality by telescoping.

Step 3: squeeze at intermediate scales. Let \(s \in (0,1]\) and choose \(n \in \mathbb{Z}^{+}\) with

\begin{equation*} \frac{1}{n+1} < s \leq \frac{1}{n} \end{equation*}

(such an \(n\) exists and is unique: take \(n = \lfloor 1/s \rfloor\), which lies in \(\mathbb{Z}^{+}\) because \(1/s \geq 1\)). Then \(\big(0,\frac{1}{n+1}\big) \subseteq (0,s) \subseteq \big(0,\frac{1}{n}\big]\), and the single point \(\frac{1}{n}\) has outer measure \(0\), so monotonicity of outer measure and Step 2 give

\begin{equation*} \frac{c}{n+1} = \Big|F \cap \big(0,\tfrac{1}{n+1}\big)\Big| \leq |F \cap (0,s)| \leq \Big|F \cap \big(0,\tfrac{1}{n}\big)\Big| = \frac{c}{n}. \end{equation*}

Dividing by \(s\) and using \(\frac{1}{n+1} < s \leq \frac{1}{n}\) (so that \(\frac{1}{s} < n+1\) and \(\frac{1}{s} \geq n\)) yields

\begin{equation*} \frac{cn}{n+1} = \frac{c/(n+1)}{1/n} \leq \frac{|F \cap (0,s)|}{s} \leq \frac{c/n}{1/(n+1)} = \frac{c(n+1)}{n}. \end{equation*}

As \(s \downarrow 0\) we have \(n = \lfloor 1/s \rfloor \to \infty\), and

\begin{equation*} \frac{cn}{n+1} \to c, \qquad \frac{c(n+1)}{n} \to c. \end{equation*}

Because \(c \geq 0\), the lower bound \(\frac{cn}{n+1}\) increases with \(n\) and the upper bound \(\frac{c(n+1)}{n}\) decreases with \(n\). Hence, given \(\varepsilon > 0\), choose \(N \in \mathbb{Z}^{+}\) with \(\frac{cN}{N+1} > c - \varepsilon\) and \(\frac{c(N+1)}{N} < c + \varepsilon\); every \(s \in \big(0,\frac{1}{N}\big]\) has \(n = \lfloor 1/s \rfloor \geq N\) and therefore satisfies \(c - \varepsilon < \frac{cn}{n+1} \leq \frac{|F \cap (0,s)|}{s} \leq \frac{c(n+1)}{n} < c + \varepsilon\). Thus

\begin{equation*} \Big| \frac{|F \cap (0,s)|}{s} - c \Big| < \varepsilon \quad \text{for all } s \in \big(0, \tfrac{1}{N}\big]. \end{equation*}

Therefore \(\lim_{s \downarrow 0} \frac{|F \cap (0,s)|}{s} = c\), and by Step 1 the density of \(E\) at \(0\) equals \(c\), as required.

Remarks. The two extreme cases are consistent with the general computation: for \(c = 0\) we get \(F = \varnothing\) and \(E = \varnothing\), whose density at \(0\) is \(0\); for \(c = 1\) we get \(F = (0,1)\) and \(E = (-1,0) \cup (0,1)\), whose density at \(0\) is \(1\) (of course \(E = \mathbb{R}\) would serve equally well). Note also that no set can have density at \(0\) outside \([0,1]\), since \(0 \leq \frac{|E \cap (-s,s)|}{2s} \leq 1\) for every \(s > 0\); so this exercise shows that every conceivable value of the density is attained.

The essential point in the construction is that the length \(\frac{c}{k(k+1)}\) of the \(k\)-th block is small compared with its distance \(\frac{1}{k+1}\) from the origin: their ratio is \(\frac{c}{k}\), which tends to \(0\). That is what makes the squeeze in Step 3 close up, since truncating in the middle of a block then changes \(\frac{|F \cap (0,s)|}{s}\) by an amount tending to \(0\). Blocks placed on a geometric scale, say occupying the fraction \(c\) of each interval \(\big[2^{-k-1}, 2^{-k}\big)\), would instead give a ratio invariant under \(s \mapsto s/2\) and not constant in \(s\), hence with no limit as \(s \downarrow 0\); that is how one answers Exercise 7.

Problem (4B.10)

Suppose \(E\) is a Lebesgue measurable subset of \(\mathbb{R}\) such that the density of \(E\) equals \(1\) at every element of \(E\) and equals \(0\) at every element of \(\mathbb{R} \setminus E\). Prove that \(E = \varnothing\) or \(E = \mathbb{R}\).

Solution

The Lebesgue Density Theorem (4.24) says that every Lebesgue measurable set \(E\) has density \(1\) at almost every element of \(E\) and density \(0\) at almost every element of \(\mathbb{R} \setminus E\). This exercise says that the exceptional sets of measure \(0\) can be removed only in the two trivial cases: a set that is its own set of density points, and whose complement is likewise, must be \(\varnothing\) or \(\mathbb{R}\).

The strategy: build from \(E\) an everywhere-differentiable function whose derivative is exactly \(\chi_E\), and then use the fact that a derivative cannot omit an intermediate value.

Step 1: an antiderivative of \(\chi_E\). Define \(g : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} g(x) = \begin{cases} \ \ \, |E \cap [0,x]| & \text{if } x \geq 0, \\[2pt] -|E \cap [x,0]| & \text{if } x < 0. \end{cases} \end{equation*}

The Lebesgue measurable sets form a \(\sigma\)-algebra on which outer measure is a measure (2.72), so each set \(E \cap [u,v]\) is Lebesgue measurable, and \(|E \cap [u,v]| \leq v - u < \infty\) by monotonicity of outer measure and 2.14. Thus \(g\) is a well-defined real-valued function. Below, every use of the word additivity means countable additivity of the measure supplied by 2.72, applied to disjoint Lebesgue measurable sets; all the sets involved have finite measure, so the subtractions performed are legitimate.

If \(u < v\), then

\begin{equation*} g(v) - g(u) = |E \cap [u,v]|. \end{equation*}

Indeed, if \(0 \leq u < v\) this is additivity applied to the disjoint measurable sets \(E \cap [0,u]\) and \(E \cap (u,v]\), whose union is \(E \cap [0,v]\), together with \(|E \cap (u,v]| = |E \cap [u,v]|\) (the single point \(u\) has measure \(0\)); if \(u < v < 0\) it is the same computation for \(E \cap [u,0] = (E \cap [u,v)) \cup (E \cap [v,0])\); and if \(u < 0 \leq v\) then \(g(v) - g(u) = |E \cap [0,v]| + |E \cap [u,0]| = |E \cap [u,v]|\), again because \(\{0\}\) is a null set.

In particular \(0 \leq g(v) - g(u) \leq v - u\) whenever \(u < v\), so \(g\) is nondecreasing and Lipschitz with constant \(1\); in particular \(g\) is continuous.

Step 2: \(g\) is differentiable at every point, with \(g^{\prime} = \chi_E\). Let \(b \in \mathbb{R}\) and let \(t > 0\). Using Step 1 twice,

\begin{equation*} \frac{g(b+t) - g(b)}{t} = \frac{|E \cap [b, b+t]|}{t}, \qquad \frac{g(b) - g(b-t)}{t} = \frac{|E \cap [b-t, b]|}{t}. \end{equation*}

Both quantities lie between \(0\) and \(1\).

First suppose \(b \in E\), so the density of \(E\) at \(b\) equals \(1\). Since \([b, b+t] \subseteq [b-t, b+t]\) and \(E\) is measurable,

\begin{equation*} 0 \leq 1 - \frac{|E \cap [b,b+t]|}{t} = \frac{\big|[b,b+t] \setminus E\big|}{t} \leq \frac{\big|[b-t,b+t] \setminus E\big|}{t} = \frac{2t - |E \cap (b-t,b+t)|}{t}, \end{equation*}

where the first equality is additivity applied to the disjoint measurable sets \(E \cap [b,b+t]\) and \([b,b+t] \setminus E\), whose union is \([b,b+t]\) (of measure \(t\) by 2.14); the middle inequality is monotonicity, since \([b,b+t] \setminus E \subseteq [b-t,b+t] \setminus E\); and the last equality is additivity applied to \(E \cap [b-t,b+t]\) and \([b-t,b+t] \setminus E\) inside \([b-t,b+t]\) (of measure \(2t\) by 2.14), together with the fact that discarding the two endpoints changes no measure, so \(|E \cap [b-t,b+t]| = |E \cap (b-t,b+t)|\). The right-hand side equals

\begin{equation*} 2\Big(1 - \frac{|E \cap (b-t,b+t)|}{2t}\Big), \end{equation*}

which tends to \(0\) as \(t \downarrow 0\) because the density of \(E\) at \(b\) is \(1\). The identical bound applies with \([b,b+t]\) replaced by \([b-t,b]\), since that interval is also contained in \([b-t,b+t]\). Hence both one-sided difference quotients displayed above tend to \(1\) as \(t \downarrow 0\); that is, \(g\) is differentiable at \(b\) and \(g^{\prime}(b) = 1\).

Now suppose \(b \in \mathbb{R} \setminus E\), so the density of \(E\) at \(b\) equals \(0\). Then

\begin{equation*} 0 \leq \frac{|E \cap [b,b+t]|}{t} = \frac{|E \cap (b,b+t)|}{t} \leq \frac{|E \cap (b-t,b+t)|}{t} = 2 \cdot \frac{|E \cap (b-t,b+t)|}{2t} \longrightarrow 0 \end{equation*}

as \(t \downarrow 0\), and the same computation applies to \(|E \cap [b-t,b]|\). Hence \(g\) is differentiable at \(b\) and \(g^{\prime}(b) = 0\).

Combining the two cases: \(g\) is differentiable at every point of \(\mathbb{R}\) and

\begin{equation*} g^{\prime}(b) = \chi_E(b) \quad \text{for every } b \in \mathbb{R}. \end{equation*}

Note that this is an assertion at every point of \(\mathbb{R}\), not merely at almost every point; that is precisely what the strong hypothesis on \(E\) buys us, and it is what makes the next step work. Compare 4.19, which for an arbitrary \(f \in L^1(\mathbb{R})\) gives only the corresponding almost-everywhere statement (and which does not apply directly here, because \(\chi_E\) need not be in \(L^1(\mathbb{R})\)).

Step 3: a derivative cannot skip the value \(\frac{1}{2}\). Suppose, for contradiction, that \(E \neq \varnothing\) and \(E \neq \mathbb{R}\). Choose \(p \in E\) and \(q \in \mathbb{R} \setminus E\).

We may assume \(p < q\). Indeed, if \(q < p\), replace \(E\) by \(-E = \{-x : x \in E\}\). Lebesgue measure is invariant under the reflection \(x \mapsto -x\), so for every \(b \in \mathbb{R}\) and \(t > 0\) we have \(|(-E) \cap (-b-t, -b+t)| = |E \cap (b-t,b+t)|\); hence the density of \(-E\) at \(-b\) equals the density of \(E\) at \(b\), and therefore \(-E\) satisfies the same hypotheses as \(E\). Moreover \(-p \in -E\) and \(-q \notin -E\) with \(-p < -q\), while \(-E = \varnothing\) exactly when \(E = \varnothing\) and \(-E = \mathbb{R}\) exactly when \(E = \mathbb{R}\). So it does no harm to assume \(p < q\).

Define \(\varphi : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} \varphi(x) = g(x) - \tfrac{1}{2} x . \end{equation*}

By Step 2, \(\varphi\) is differentiable at every point, with \(\varphi^{\prime}(x) = \chi_E(x) - \frac{1}{2}\). In particular

\begin{equation*} \varphi^{\prime}(p) = \tfrac{1}{2} > 0 \qquad \text{and} \qquad \varphi^{\prime}(q) = -\tfrac{1}{2} < 0 . \end{equation*}

Since \(g\) is continuous, so is \(\varphi\); as \([p,q]\) is compact, \(\varphi\) attains a maximum on \([p,q]\), say at \(c \in [p,q]\).

The maximum is not attained at \(p\): because

\begin{equation*} \lim_{t \downarrow 0} \frac{\varphi(p+t) - \varphi(p)}{t} = \varphi^{\prime}(p) = \tfrac{1}{2} > 0, \end{equation*}

there exists \(t \in (0, q-p)\) with \(\frac{\varphi(p+t) - \varphi(p)}{t} > 0\), hence \(\varphi(p+t) > \varphi(p)\) with \(p + t \in (p,q)\).

The maximum is not attained at \(q\): because

\begin{equation*} \lim_{t \uparrow 0} \frac{\varphi(q+t) - \varphi(q)}{t} = \varphi^{\prime}(q) = -\tfrac{1}{2} < 0, \end{equation*}

there exists \(t \in (p-q, 0)\) with \(\frac{\varphi(q+t) - \varphi(q)}{t} < 0\); since \(t < 0\), this gives \(\varphi(q+t) > \varphi(q)\) with \(q + t \in (p,q)\).

Hence \(c \in (p,q)\), an interior maximum point. Since \(\varphi\) is differentiable at \(c\) we get, considering the two one-sided difference quotients,

\begin{equation*} \varphi^{\prime}( c) = \lim_{t \downarrow 0} \frac{\varphi(c+t) - \varphi( c)}{t} \leq 0 \qquad \text{and} \qquad \varphi^{\prime}( c) = \lim_{t \uparrow 0} \frac{\varphi(c+t) - \varphi( c)}{t} \geq 0, \end{equation*}

because the numerators are \(\leq 0\) for all \(t\) with \(c + t \in [p,q]\) (as \(c\) is a maximum point of \(\varphi\) on \([p,q]\), and \(c\) is interior so both signs of small \(t\) are allowed), while \(t > 0\) in the first limit and \(t < 0\) in the second. Therefore \(\varphi^{\prime}( c) = 0\), which means

\begin{equation*} g^{\prime}( c) = \tfrac{1}{2}. \end{equation*}

But Step 2 says \(g^{\prime}( c) = \chi_E( c) \in \{0, 1\}\), a contradiction.

Consequently no such pair \(p, q\) exists, which means \(E = \varnothing\) or \(E = \mathbb{R}\).

Remarks. (a) Step 3 is just the proof of Darboux’s theorem (derivatives have the intermediate value property) specialized to the value \(\frac{1}{2}\); the point is that \(g^{\prime}\) takes the values \(1\) (at \(p\)) and \(0\) (at \(q\)) but, by Step 2, never takes any value strictly between them.

(b) A soft version of the conclusion is already visible without Darboux’s theorem: if \(|E| = 0\), then \(\frac{|E \cap (b-t,b+t)|}{2t} = 0\) for all \(b\) and all \(t > 0\), so the density of \(E\) is \(0\) at every point, so no point can be an element of \(E\) (each element of \(E\) must have density \(1\)), giving \(E = \varnothing\). Symmetrically, if \(|\mathbb{R} \setminus E| = 0\) then the density of \(E\) is \(1\) everywhere, forcing \(\mathbb{R} \setminus E = \varnothing\). The content of the exercise is to rule out the remaining possibility, that \(E\) and its complement both have positive measure, and for that the differentiability argument above does the work.

(c) In the language of point-set topology, the measurable sets each of whose points is a density point of the set form a topology on \(\mathbb{R}\), called the density topology; it is finer than the usual topology, and this exercise says that \(\mathbb{R}\) is connected in it.

Product Measures

Exercises 5A

Problem (5A.1)

Suppose \((X, \mathcal{S})\) and \((Y, \mathcal{T})\) are measurable spaces. Prove that if \(A\) is a nonempty subset of \(X\) and \(B\) is a nonempty subset of \(Y\) such that \(A \times B \in \mathcal{S} \otimes \mathcal{T}\), then \(A \in \mathcal{S}\) and \(B \in \mathcal{T}\).

Solution

Suppose \(A \subseteq X\) and \(B \subseteq Y\) are nonempty and \(A \times B \in \mathcal{S} \otimes \mathcal{T}\).

Because \(B\) is nonempty, we may choose \(b \in B\). By 5.6, the cross section \([A \times B]^b\) belongs to \(\mathcal{S}\) (this is where the hypothesis \(A \times B \in \mathcal{S} \otimes \mathcal{T}\) is used). On the other hand, Example 5.5 tells us that

\begin{equation*} [A \times B]^b = \begin{cases} A & \text{if } b \in B,\\ \varnothing & \text{if } b \notin B.\end{cases} \end{equation*}

Since \(b \in B\), we have \([A \times B]^b = A\). Hence \(A \in \mathcal{S}\).

For the reader who prefers to see the cross section computed directly: \(x \in [A \times B]^b\) means \((x, b) \in A \times B\), which holds if and only if \(x \in A\) and \(b \in B\); as \(b \in B\), this is equivalent to \(x \in A\).

Similarly, because \(A\) is nonempty we may choose \(a \in A\). By 5.6 we have \([A \times B]_a \in \mathcal{T}\), and by Example 5.5,

\begin{equation*} [A \times B]_a = \begin{cases} B & \text{if } a \in A,\\ \varnothing & \text{if } a \notin A,\end{cases} \end{equation*}

so \([A \times B]_a = B\). Hence \(B \in \mathcal{T}\).

Thus \(A \in \mathcal{S}\) and \(B \in \mathcal{T}\), as desired.

The hypothesis that \(A\) and \(B\) are nonempty cannot be dropped: if \(A\) is an arbitrary (possibly non-measurable) subset of \(X\) and \(B = \varnothing\), then \(A \times B = \varnothing \in \mathcal{S} \otimes \mathcal{T}\), even though \(A\) need not be in \(\mathcal{S}\).

Problem (5A.2)

Suppose \((X, \mathcal{S})\) is a measurable space. Prove that if \(E \in \mathcal{S} \otimes \mathcal{S}\), then

\begin{equation*} \{x \in X : (x, x) \in E\} \in \mathcal{S}. \end{equation*}

Solution

For \(E \subseteq X \times X\) write

\begin{equation*} \Delta(E) = \{x \in X : (x,x) \in E\}. \end{equation*}

We use the standard two-step technique described just before 5.10: we show that all measurable rectangles have the desired property, and that the collection of sets with the desired property is a \(\sigma\)-algebra.

Let

\begin{equation*} \mathcal{E} = \{E \subseteq X \times X : \Delta(E) \in \mathcal{S}\}. \end{equation*}

Measurable rectangles belong to \(\mathcal{E}\). Indeed, suppose \(A, B \in \mathcal{S}\). Then for \(x \in X\) we have \((x,x) \in A \times B\) if and only if \(x \in A\) and \(x \in B\). Hence

\begin{equation*} \Delta(A \times B) = A \cap B, \end{equation*}

which is in \(\mathcal{S}\) because \(\sigma\)-algebras are closed under finite intersections. Thus \(A \times B \in \mathcal{E}\).

\(\mathcal{E}\) contains \(\varnothing\), since \(\Delta(\varnothing) = \varnothing \in \mathcal{S}\).

\(\mathcal{E}\) is closed under complementation. Suppose \(E \in \mathcal{E}\). For \(x \in X\) we have \((x,x) \in (X \times X) \setminus E\) if and only if \((x,x) \notin E\). Hence

\begin{equation*} \Delta\bigl((X \times X) \setminus E\bigr) = X \setminus \Delta(E), \end{equation*}

which is in \(\mathcal{S}\) because \(\Delta(E) \in \mathcal{S}\) and \(\mathcal{S}\) is closed under complementation. Thus \((X \times X) \setminus E \in \mathcal{E}\).

\(\mathcal{E}\) is closed under countable unions. Suppose \(E_1, E_2, \ldots \in \mathcal{E}\). For \(x \in X\) we have \((x,x) \in \bigcup_{k=1}^{\infty} E_k\) if and only if \((x,x) \in E_k\) for some \(k\). Hence

\begin{equation*} \Delta\Bigl(\bigcup_{k=1}^{\infty} E_k\Bigr) = \bigcup_{k=1}^{\infty} \Delta(E_k), \end{equation*}

which is in \(\mathcal{S}\) because each \(\Delta(E_k) \in \mathcal{S}\) and \(\mathcal{S}\) is closed under countable unions. Thus \(\bigcup_{k=1}^{\infty} E_k \in \mathcal{E}\).

The three previous paragraphs show that \(\mathcal{E}\) is a \(\sigma\)-algebra on \(X \times X\), and the paragraph before them shows that \(\mathcal{E}\) contains every measurable rectangle in \(\mathcal{S} \otimes \mathcal{S}\). Because \(\mathcal{S} \otimes \mathcal{S}\) is by definition the smallest \(\sigma\)-algebra on \(X \times X\) containing all the measurable rectangles, we conclude that

\begin{equation*} \mathcal{S} \otimes \mathcal{S} \subseteq \mathcal{E}. \end{equation*}

In other words, if \(E \in \mathcal{S} \otimes \mathcal{S}\), then \(\{x \in X : (x,x) \in E\} \in \mathcal{S}\), as desired.

Equivalently, one can phrase the argument as follows: the map \(\varphi : X \to X \times X\) defined by \(\varphi(x) = (x,x)\) satisfies \(\varphi^{-1}(A \times B) = A \cap B \in \mathcal{S}\) for all \(A, B \in \mathcal{S}\), and \(\{E \subseteq X \times X : \varphi^{-1}(E) \in \mathcal{S}\}\) is a \(\sigma\)-algebra (because taking inverse images commutes with complementation and unions); hence \(\varphi^{-1}(E) \in \mathcal{S}\) for every \(E \in \mathcal{S} \otimes \mathcal{S}\), and \(\varphi^{-1}(E) = \Delta(E)\).

Problem (5A.3)

Let \(\mathcal{B}\) denote the \(\sigma\)-algebra of Borel subsets of \(\mathbb{R}\). Show that there exists a set \(E \subseteq \mathbb{R} \times \mathbb{R}\) such that \([E]_a \in \mathcal{B}\) and \([E]^a \in \mathcal{B}\) for every \(a \in \mathbb{R}\), but \(E \notin \mathcal{B} \otimes \mathcal{B}\).

Solution

By 2.67 there exists a set \(A \subseteq \mathbb{R}\) that is not a Borel set. Fix such an \(A\) and define

\begin{equation*} E = \{(x, x) : x \in A\}, \end{equation*}

so that \(E\) is the part of the diagonal of \(\mathbb{R} \times \mathbb{R}\) sitting over \(A\).

First we compute the cross sections. Let \(a \in \mathbb{R}\). By definition (5.3),

\begin{equation*} [E]_a = \{y \in \mathbb{R} : (a, y) \in E\}. \end{equation*}

Now \((a,y) \in E\) means \((a,y) = (x,x)\) for some \(x \in A\), which forces \(x = a = y\) and \(a \in A\). Hence

\begin{equation*} [E]_a = \begin{cases} \{a\} & \text{if } a \in A,\\ \varnothing & \text{if } a \notin A.\end{cases} \end{equation*}

In either case \([E]_a\) is a closed subset of \(\mathbb{R}\) and therefore \([E]_a \in \mathcal{B}\).

Because \(E\) is a subset of the diagonal, it is symmetric in the two coordinates in the following sense: \((x,y) \in E\) if and only if \((y,x) \in E\). Consequently the same computation gives

\begin{equation*} [E]^a = \{x \in \mathbb{R} : (x, a) \in E\} = \begin{cases} \{a\} & \text{if } a \in A,\\ \varnothing & \text{if } a \notin A,\end{cases} \end{equation*}

so \([E]^a \in \mathcal{B}\) as well. Thus every cross section of \(E\), of either kind, is a Borel subset of \(\mathbb{R}\).

Now suppose, seeking a contradiction, that \(E \in \mathcal{B} \otimes \mathcal{B}\). Apply Exercise 5A.2 with \(X = \mathbb{R}\) and \(\mathcal{S} = \mathcal{B}\): that exercise states that if \(E \in \mathcal{B} \otimes \mathcal{B}\), then

\begin{equation*} \{x \in \mathbb{R} : (x,x) \in E\} \in \mathcal{B}. \end{equation*}

But by the definition of \(E\) we have \((x,x) \in E\) if and only if \(x \in A\), and hence

\begin{equation*} \{x \in \mathbb{R} : (x,x) \in E\} = A. \end{equation*}

Thus \(A \in \mathcal{B}\), contradicting the choice of \(A\) as a non-Borel set. Therefore \(E \notin \mathcal{B} \otimes \mathcal{B}\).

This \(E\) has all cross sections Borel yet is not in \(\mathcal{B} \otimes \mathcal{B}\), which is what we wanted. In particular the converse of 5.6 is false: measurability of all the cross sections of a set does not imply measurability of the set in the product \(\sigma\)-algebra.

Problem (5A.4)

Suppose \((X, \mathcal{S})\) and \((Y, \mathcal{T})\) are measurable spaces. Prove that if \(f : X \to \mathbb{R}\) is \(\mathcal{S}\)-measurable and \(g : Y \to \mathbb{R}\) is \(\mathcal{T}\)-measurable and \(h : X \times Y \to \mathbb{R}\) is defined by \(h(x,y) = f(x)\, g(y)\), then \(h\) is \((\mathcal{S} \otimes \mathcal{T})\)-measurable.

Solution

Define \(F, G : X \times Y \to \mathbb{R}\) by

\begin{equation*} F(x,y) = f(x) \quad\text{and}\quad G(x,y) = g(y) \end{equation*}

for \((x,y) \in X \times Y\). Then \(h = F G\), so by 2.46(a) (products of \(\mathcal{S} \otimes \mathcal{T}\)-measurable functions are \(\mathcal{S} \otimes \mathcal{T}\)-measurable) it suffices to prove that \(F\) and \(G\) are \((\mathcal{S} \otimes \mathcal{T})\)-measurable.

Suppose \(D \subseteq \mathbb{R}\) is a Borel set. Then

\begin{equation*} F^{-1}(D) = \{(x,y) \in X \times Y : f(x) \in D\} = f^{-1}(D) \times Y. \end{equation*}

To verify this equality: \((x,y) \in F^{-1}(D)\) if and only if \(F(x,y) = f(x) \in D\), which holds if and only if \(x \in f^{-1}(D)\) (with \(y \in Y\) arbitrary). Because \(f\) is \(\mathcal{S}\)-measurable, \(f^{-1}(D) \in \mathcal{S}\); also \(Y \in \mathcal{T}\). Hence \(f^{-1}(D) \times Y\) is a measurable rectangle in \(\mathcal{S} \otimes \mathcal{T}\) and thus belongs to \(\mathcal{S} \otimes \mathcal{T}\) (by 5.2). Because \(D\) was an arbitrary Borel subset of \(\mathbb{R}\), we conclude that \(F\) is \((\mathcal{S} \otimes \mathcal{T})\)-measurable.

Similarly, for a Borel set \(D \subseteq \mathbb{R}\),

\begin{equation*} G^{-1}(D) = \{(x,y) \in X \times Y : g(y) \in D\} = X \times g^{-1}(D), \end{equation*}

which is a measurable rectangle in \(\mathcal{S} \otimes \mathcal{T}\) because \(X \in \mathcal{S}\) and \(g^{-1}(D) \in \mathcal{T}\) (the latter because \(g\) is \(\mathcal{T}\)-measurable). Hence \(G\) is \((\mathcal{S} \otimes \mathcal{T})\)-measurable.

Now 2.46(a), applied to the measurable space \((X \times Y, \mathcal{S} \otimes \mathcal{T})\) and the two \((\mathcal{S} \otimes \mathcal{T})\)-measurable functions \(F\) and \(G\), shows that \(FG\) is \((\mathcal{S} \otimes \mathcal{T})\)-measurable. Because

\begin{equation*} (FG)(x,y) = F(x,y)\, G(x,y) = f(x)\, g(y) = h(x,y) \end{equation*}

for all \((x,y) \in X \times Y\), we have \(h = FG\). Hence \(h\) is \((\mathcal{S} \otimes \mathcal{T})\)-measurable, as desired.

Problem (5A.5)

Verify the assertion in Example 5.11 that the collection of finite unions of intervals of \(\mathbb{R}\) is closed under complementation.

Solution

Let \(\mathcal{A}\) denote the collection of all finite unions of intervals of \(\mathbb{R}\), where as in Example 5.11 the word interval includes open, closed, half-open, bounded, and unbounded intervals, as well as sets consisting of a single point and the empty set (note that \(\varnothing = (0,0)\) and \(\mathbb{R} = (-\infty, \infty)\) are intervals).

Call a set \(I \subseteq \mathbb{R}\) order convex if

\begin{equation*} x, z \in I \text{ and } x < y < z \implies y \in I. \end{equation*}

Step 0 below shows that the order convex subsets of \(\mathbb{R}\) are exactly the intervals; this characterization is what makes the rest of the proof work, so we prove it rather than assume it.

Step 0: a set \(I \subseteq \mathbb{R}\) is an interval if and only if \(I\) is order convex.

Every interval is order convex: each of the forms \((\alpha, \beta)\), \([\alpha, \beta]\), \([\alpha, \beta)\), \((\alpha, \beta]\) (allowing \(\alpha = -\infty\) and \(\beta = \infty\) for the open ends), along with singletons and \(\varnothing\), is defined by a pair of inequalities that any \(y\) strictly between two of its elements also satisfies.

Conversely, suppose \(I\) is order convex. If \(I = \varnothing\) then \(I = (0,0)\) is an interval. So assume \(I \ne \varnothing\) and set

\begin{equation*} \alpha = \inf I \in [-\infty, \infty), \qquad \beta = \sup I \in (-\infty, \infty], \end{equation*}

where \(\alpha = -\infty\) if \(I\) is unbounded below and \(\beta = \infty\) if \(I\) is unbounded above. We claim

\begin{equation*} (\alpha, \beta) \subseteq I \subseteq [\alpha, \beta] \cap \mathbb{R}. \end{equation*}

The right inclusion is immediate from the definitions of infimum and supremum. For the left inclusion, suppose \(\alpha < y < \beta\). Because \(y > \alpha = \inf I\), the number \(y\) is not a lower bound for \(I\), so there exists \(x \in I\) with \(x < y\). Because \(y < \beta = \sup I\), the number \(y\) is not an upper bound for \(I\), so there exists \(z \in I\) with \(y < z\). Now \(x, z \in I\) and \(x < y < z\), so order convexity gives \(y \in I\).

Hence \(I\) is obtained from \((\alpha, \beta)\) by possibly adjoining the finite endpoints among \(\alpha\) and \(\beta\) that happen to lie in \(I\), which makes \(I\) equal to \((\alpha, \beta)\), \([\alpha, \beta)\), \((\alpha, \beta]\), or \([\alpha, \beta]\) (an infinite endpoint is never adjoined, as \(I \subseteq \mathbb{R}\)). Each of these is an interval in the sense of Example 5.11; note that when \(\alpha = \beta\) the only possibility is \(I = [\alpha, \alpha] = \{\alpha\}\), a single point, which Example 5.11 explicitly counts as an interval. This proves Step 0.

From now on we use interval and order convex interchangeably, as Step 0 permits.

Step 1: \(\mathcal{A}\) is closed under finite intersections.

If \(I\) and \(J\) are intervals, then \(I \cap J\) is an interval: by Step 0 it is enough to check that \(I \cap J\) is order convex, so suppose \(x, z \in I \cap J\) and \(x < y < z\); then \(y \in I\) (because \(x, z \in I\) and \(I\) is order convex) and \(y \in J\) (because \(x, z \in J\) and \(J\) is order convex), so \(y \in I \cap J\). Now if \(I_1, \ldots, I_n\) and \(J_1, \ldots, J_m\) are intervals, then

\begin{equation*} \Bigl(\bigcup_{j=1}^{n} I_j\Bigr) \cap \Bigl(\bigcup_{k=1}^{m} J_k\Bigr) = \bigcup_{j=1}^{n} \bigcup_{k=1}^{m} (I_j \cap J_k), \end{equation*}

which is a finite union of intervals and therefore is in \(\mathcal{A}\). Thus \(\mathcal{A}\) is closed under the intersection of two of its elements, and hence, by induction, under the intersection of finitely many of its elements.

Step 2: the complement of an interval is the union of two intervals.

Let \(I\) be an interval and define

\begin{equation*} L = \{x \in \mathbb{R} : x < y \text{ for every } y \in I\}, \qquad U = \{x \in \mathbb{R} : x > y \text{ for every } y \in I\}. \end{equation*}

We claim \(\mathbb{R} \setminus I = L \cup U\).

For the inclusion \(L \cup U \subseteq \mathbb{R} \setminus I\): if \(x \in L\) and we had \(x \in I\), then taking \(y = x\) in the definition of \(L\) would give \(x < x\), a contradiction; hence \(x \notin I\). The same argument applies to \(U\).

For the inclusion \(\mathbb{R} \setminus I \subseteq L \cup U\): suppose \(x \notin I\) but \(x \notin L\) and \(x \notin U\). From \(x \notin L\) we get some \(y \in I\) with \(y \le x\); since \(x \notin I\) we have \(y \ne x\), so \(y < x\). From \(x \notin U\) we get some \(z \in I\) with \(z \ge x\), and again \(z \ne x\), so \(x < z\). Thus \(y, z \in I\) and \(y < x < z\), which by the order convexity of \(I\) forces \(x \in I\), a contradiction. Hence \(x \in L \cup U\).

Finally, \(L\) and \(U\) are intervals. The set \(L\) is closed downward: if \(x < x^{\prime}\) and \(x^{\prime} \in L\), then \(x < x^{\prime} < y\) for every \(y \in I\), so \(x \in L\). A downward closed set is order convex, because if \(x, z \in L\) and \(x < y < z\), then \(y < z\) and \(z \in L\) give \(y \in L\); hence \(L\) is an interval by Step 0. Symmetrically, \(U\) is closed upward, hence order convex, hence an interval. (Concretely, if \(I\) is nonempty then \(L\) is \((-\infty, \inf I)\) or \((-\infty, \inf I]\) and \(U\) is \((\sup I, \infty)\) or \([\sup I, \infty)\), where these are interpreted as \(\varnothing\) when \(\inf I = -\infty\) or \(\sup I = \infty\); if \(I = \varnothing\) then \(L = U = \mathbb{R}\).)

Step 3: \(\mathcal{A}\) is closed under complementation.

Suppose \(E \in \mathcal{A}\), say \(E = I_1 \cup \cdots \cup I_n\) with each \(I_j\) an interval. By De Morgan’s Laws,

\begin{equation*} \mathbb{R} \setminus E = \bigcap_{j=1}^{n} (\mathbb{R} \setminus I_j). \end{equation*}

By Step 2, each \(\mathbb{R} \setminus I_j\) is a union of two intervals and hence is in \(\mathcal{A}\). By Step 1, the intersection of these finitely many elements of \(\mathcal{A}\) is in \(\mathcal{A}\). Hence \(\mathbb{R} \setminus E \in \mathcal{A}\).

Because \(\mathcal{A}\) is obviously closed under finite unions and contains \(\varnothing\), and we have now shown that it is closed under complementation, \(\mathcal{A}\) is an algebra on \(\mathbb{R}\), as asserted in Example 5.11.

Problem (5A.6)

Verify the assertion in Example 5.12 that the collection of countable unions of intervals of \(\mathbb{R}\) is not closed under complementation.

Solution

Let \(\mathcal{A}\) denote the collection of all countable unions of intervals of \(\mathbb{R}\) (as in Example 5.11, single-point sets count as intervals). We exhibit a set in \(\mathcal{A}\) whose complement is not in \(\mathcal{A}\).

The set \(\mathbb{Q}\) of rational numbers is in \(\mathcal{A}\). Indeed, \(\mathbb{Q}\) is countable, and

\begin{equation*} \mathbb{Q} = \bigcup_{q \in \mathbb{Q}} \{q\}, \end{equation*}

which exhibits \(\mathbb{Q}\) as a countable union of the intervals \(\{q\}\) (each a set consisting of a single point).

We claim that \(\mathbb{R} \setminus \mathbb{Q}\), the set of irrational numbers, is not in \(\mathcal{A}\).

First observe that if \(I\) is an interval with \(I \subseteq \mathbb{R} \setminus \mathbb{Q}\), then \(I\) contains at most one point. To see this, suppose \(x, z \in I\) with \(x < z\). Because \(\mathbb{Q}\) is dense in \(\mathbb{R}\), there exists \(q \in \mathbb{Q}\) with \(x < q < z\). Every interval \(I\) has the property that \(x, z \in I\) and \(x < y < z\) imply \(y \in I\), because each of the possible forms \((\alpha, \beta)\), \([\alpha, \beta)\), \((\alpha, \beta]\), \([\alpha, \beta]\) (with \(\alpha = -\infty\) or \(\beta = \infty\) allowed at an open end), together with singletons and \(\varnothing\), is cut out by inequalities that \(y\) inherits from \(x\) and \(z\). Applying this with \(y = q\) gives \(q \in I\), contradicting \(I \cap \mathbb{Q} = \varnothing\). Hence \(I\) cannot contain two distinct points.

Now suppose, seeking a contradiction, that \(\mathbb{R} \setminus \mathbb{Q} \in \mathcal{A}\), say

\begin{equation*} \mathbb{R} \setminus \mathbb{Q} = \bigcup_{k=1}^{\infty} I_k \end{equation*}

with each \(I_k\) an interval (a countable union can always be written as a union indexed by \(\mathbb{Z}^{+}\), repeating terms or using \(\varnothing\) if the index set is finite). Each \(I_k\) satisfies \(I_k \subseteq \mathbb{R} \setminus \mathbb{Q}\), so by the previous paragraph each \(I_k\) has at most one element. Therefore \(\bigcup_{k=1}^{\infty} I_k\) is a countable union of sets each having at most one element, and hence is a countable set.

But \(\mathbb{R} \setminus \mathbb{Q}\) is uncountable: if it were countable, then \(\mathbb{R} = \mathbb{Q} \cup (\mathbb{R} \setminus \mathbb{Q})\) would be the union of two countable sets and thus countable, contradicting the uncountability of \(\mathbb{R}\).

This contradiction shows that \(\mathbb{R} \setminus \mathbb{Q} \notin \mathcal{A}\). Since \(\mathbb{Q} \in \mathcal{A}\), the collection \(\mathcal{A}\) is not closed under complementation, as asserted in Example 5.12. In particular \(\mathcal{A}\) is neither an algebra nor a \(\sigma\)-algebra on \(\mathbb{R}\), even though it is closed under countable unions.

Problem (5A.7)

Suppose \(\mathcal{A}\) is a nonempty collection of subsets of a set \(W\). Show that \(\mathcal{A}\) is an algebra on \(W\) if and only if \(\mathcal{A}\) is closed under finite intersections and under complementation.

Solution

Recall from 5.10 that \(\mathcal{A}\) is an algebra on \(W\) if \(\varnothing \in \mathcal{A}\), and \(W \setminus E \in \mathcal{A}\) whenever \(E \in \mathcal{A}\), and \(E \cup F \in \mathcal{A}\) whenever \(E, F \in \mathcal{A}\). As usual, closure under finite intersections means that \(E_1 \cap \cdots \cap E_n \in \mathcal{A}\) whenever \(n \in \mathbb{Z}^{+}\) and \(E_1, \ldots, E_n \in \mathcal{A}\); by induction this is equivalent to closure under the intersection of two elements of \(\mathcal{A}\).

First suppose \(\mathcal{A}\) is an algebra on \(W\). Then \(\mathcal{A}\) is closed under complementation, by the second bullet point of 5.10. To see that \(\mathcal{A}\) is closed under finite intersections, suppose \(E, F \in \mathcal{A}\). Then \(W \setminus E \in \mathcal{A}\) and \(W \setminus F \in \mathcal{A}\) (closure under complementation), hence \((W \setminus E) \cup (W \setminus F) \in \mathcal{A}\) (closure under finite unions), hence, using De Morgan’s Laws,

\begin{equation*} E \cap F = W \setminus \bigl((W \setminus E) \cup (W \setminus F)\bigr) \in \mathcal{A} \end{equation*}

(closure under complementation again). Induction now gives \(E_1 \cap \cdots \cap E_n \in \mathcal{A}\) for all \(E_1, \ldots, E_n \in \mathcal{A}\): the case \(n = 1\) is trivial, and if \(E_1 \cap \cdots \cap E_{n-1} \in \mathcal{A}\) then intersecting this set with \(E_n\) stays in \(\mathcal{A}\) by what we just proved. Thus \(\mathcal{A}\) is closed under finite intersections and under complementation. (Note that this direction does not use the hypothesis that \(\mathcal{A}\) is nonempty, which is automatic here because \(\varnothing \in \mathcal{A}\).)

Conversely, suppose \(\mathcal{A}\) is nonempty and closed under finite intersections and under complementation. We verify the three conditions in 5.10.

Because \(\mathcal{A}\) is nonempty, there exists some \(E \in \mathcal{A}\). Closure under complementation gives \(W \setminus E \in \mathcal{A}\), and then closure under finite intersections gives

\begin{equation*} \varnothing = E \cap (W \setminus E) \in \mathcal{A}. \end{equation*}

Thus the first condition of 5.10 holds. (This is exactly where the hypothesis that \(\mathcal{A}\) is nonempty is needed: the empty collection is closed under complementation and under finite intersections but is not an algebra.)

The second condition of 5.10, closure under complementation, holds by hypothesis.

For the third condition, suppose \(E, F \in \mathcal{A}\). Then \(W \setminus E \in \mathcal{A}\) and \(W \setminus F \in \mathcal{A}\), so \((W \setminus E) \cap (W \setminus F) \in \mathcal{A}\), so by De Morgan’s Laws

\begin{equation*} E \cup F = W \setminus \bigl((W \setminus E) \cap (W \setminus F)\bigr) \in \mathcal{A}. \end{equation*}

Thus \(\mathcal{A}\) is closed under the union of two of its elements, and hence (by induction, as above) under finite unions.

The three conditions of 5.10 are satisfied, so \(\mathcal{A}\) is an algebra on \(W\). This completes the proof of both implications.

Problem (5A.8)

Suppose \(\mu\) is a measure on a measurable space \((X, \mathcal{S})\). Prove that the following are equivalent.

(a) The measure \(\mu\) is \(\sigma\)-finite.

(b) There exists an increasing sequence \(X_1 \subseteq X_2 \subseteq \cdots\) of sets in \(\mathcal{S}\) such that \(X = \bigcup_{k=1}^{\infty} X_k\) and \(\mu(X_k) < \infty\) for every \(k \in \mathbf{Z}^+\).

(c) There exists a disjoint sequence \(X_1, X_2, X_3, \ldots\) of sets in \(\mathcal{S}\) such that \(X = \bigcup_{k=1}^{\infty} X_k\) and \(\mu(X_k) < \infty\) for every \(k \in \mathbf{Z}^+\).

Solution

We prove (a) \(\Rightarrow\) (b) \(\Rightarrow\) (c) \(\Rightarrow\) (a).

(a) \(\Rightarrow\) (b). Suppose \(\mu\) is \(\sigma\)-finite. By the definition of \(\sigma\)-finite (5.18), there is a sequence \(E_1, E_2, \ldots\) of sets in \(\mathcal{S}\) with

\begin{equation*} X = \bigcup_{k=1}^{\infty} E_k \quad\text{and}\quad \mu(E_k) < \infty \text{ for every } k \in \mathbf{Z}^+ . \end{equation*}

For \(k \in \mathbf{Z}^+\) define

\begin{equation*} X_k = E_1 \cup E_2 \cup \cdots \cup E_k . \end{equation*}

Each \(X_k\) is in \(\mathcal{S}\) because a \(\sigma\)-algebra is closed under finite unions, and clearly \(X_1 \subseteq X_2 \subseteq \cdots\) is an increasing sequence. Furthermore

\begin{equation*} \bigcup_{k=1}^{\infty} X_k = \bigcup_{k=1}^{\infty} E_k = X , \end{equation*}

because each \(E_k \subseteq X_k\) and each \(X_k \subseteq \bigcup_{j=1}^{\infty} E_j\). Finally, finite subadditivity of a measure (a consequence of countable subadditivity 2.58, applied to the sequence \(E_1, \ldots, E_k, \varnothing, \varnothing, \ldots\)) gives

\begin{equation*} \mu(X_k) \le \mu(E_1) + \cdots + \mu(E_k) < \infty \end{equation*}

for every \(k \in \mathbf{Z}^+\), since each of the finitely many summands is finite. Thus (b) holds.

(b) \(\Rightarrow\) (c). Suppose \(X_1 \subseteq X_2 \subseteq \cdots\) is an increasing sequence in \(\mathcal{S}\) with \(X = \bigcup_{k=1}^{\infty} X_k\) and \(\mu(X_k) < \infty\) for each \(k\). Define

\begin{equation*} Y_1 = X_1 \quad\text{and}\quad Y_k = X_k \setminus X_{k-1} \text{ for } k \ge 2 . \end{equation*}

Each \(Y_k\) is in \(\mathcal{S}\), because a \(\sigma\)-algebra is closed under complementation and finite intersections and \(X_k \setminus X_{k-1} = X_k \cap (X \setminus X_{k-1})\).

The sequence \(Y_1, Y_2, \ldots\) is disjoint: if \(j < k\), then \(Y_j \subseteq X_j \subseteq X_{k-1}\) (using that the sequence \(X_1 \subseteq X_2 \subseteq \cdots\) is increasing and \(j \le k-1\)), while \(Y_k \cap X_{k-1} = \varnothing\); hence \(Y_j \cap Y_k = \varnothing\).

The union of the \(Y_k\) is \(X\): certainly \(\bigcup_{k=1}^{\infty} Y_k \subseteq \bigcup_{k=1}^{\infty} X_k = X\). Conversely, if \(x \in X\), then \(x \in X_k\) for some \(k \in \mathbf{Z}^+\); taking \(m\) to be the smallest such \(k\), we have \(x \in X_m\) and (if \(m \ge 2\)) \(x \notin X_{m-1}\), so \(x \in Y_m\). Thus \(X = \bigcup_{k=1}^{\infty} Y_k\).

Finally \(Y_k \subseteq X_k\), so \(\mu(Y_k) \le \mu(X_k) < \infty\) by 2.57(a). Thus the disjoint sequence \(Y_1, Y_2, \ldots\) witnesses (c).

(c) \(\Rightarrow\) (a). Suppose \(X_1, X_2, \ldots\) is a disjoint sequence of sets in \(\mathcal{S}\) with \(X = \bigcup_{k=1}^{\infty} X_k\) and \(\mu(X_k) < \infty\) for every \(k \in \mathbf{Z}^+\). Then in particular \(X_1, X_2, \ldots\) is a sequence of sets in \(\mathcal{S}\) whose union is \(X\) and each of which has finite measure; disjointness is simply not used. This is exactly the definition (5.18) of \(\mu\) being \(\sigma\)-finite. Thus (a) holds.

Having shown (a) \(\Rightarrow\) (b) \(\Rightarrow\) (c) \(\Rightarrow\) (a), we conclude that (a), (b), and (c) are equivalent.

Problem (5A.9)

Suppose \(\mu\) and \(\nu\) are \(\sigma\)-finite measures. Prove that \(\mu \times \nu\) is a \(\sigma\)-finite measure.

Solution

Say \(\mu\) is a measure on the measurable space \((X, \mathcal{S})\) and \(\nu\) is a measure on the measurable space \((Y, \mathcal{T})\), both \(\sigma\)-finite; then \(\mu \times \nu\) is defined on \(\mathcal{S} \otimes \mathcal{T}\) by 5.25, and it is a measure by 5.27. So only \(\sigma\)-finiteness remains to be proved.

Because \(\mu\) is \(\sigma\)-finite, there is a sequence \(A_1, A_2, \ldots\) of sets in \(\mathcal{S}\) with

\begin{equation*} X = \bigcup_{j=1}^{\infty} A_j \quad\text{and}\quad \mu(A_j) < \infty \text{ for every } j \in \mathbf{Z}^+ . \end{equation*}

Because \(\nu\) is \(\sigma\)-finite, there is a sequence \(B_1, B_2, \ldots\) of sets in \(\mathcal{T}\) with

\begin{equation*} Y = \bigcup_{k=1}^{\infty} B_k \quad\text{and}\quad \nu(B_k) < \infty \text{ for every } k \in \mathbf{Z}^+ . \end{equation*}

Consider the family of measurable rectangles

\begin{equation*} \{ A_j \times B_k : j, k \in \mathbf{Z}^+ \} . \end{equation*}

Each \(A_j \times B_k\) belongs to \(\mathcal{S} \otimes \mathcal{T}\), being a measurable rectangle (5.2). The family is indexed by \(\mathbf{Z}^+ \times \mathbf{Z}^+\), which is a countable set, so we may list these sets as a single sequence \(E_1, E_2, \ldots\) of sets in \(\mathcal{S} \otimes \mathcal{T}\).

The union of this family is all of \(X \times Y\). Indeed, each \(A_j \times B_k \subseteq X \times Y\); conversely, if \((x, y) \in X \times Y\), then \(x \in A_j\) for some \(j \in \mathbf{Z}^+\) and \(y \in B_k\) for some \(k \in \mathbf{Z}^+\), so \((x, y) \in A_j \times B_k\). Hence

\begin{equation*} X \times Y = \bigcup_{j=1}^{\infty} \bigcup_{k=1}^{\infty} (A_j \times B_k) . \end{equation*}

Finally, by 5.26 (product measure of a measurable rectangle is the product of the measures),

\begin{equation*} (\mu \times \nu)(A_j \times B_k) = \mu(A_j)\, \nu(B_k) < \infty \end{equation*}

for all \(j, k \in \mathbf{Z}^+\), because both factors are finite real numbers.

Thus \(X \times Y\) is the union of the countable family \(E_1, E_2, \ldots\) of sets in \(\mathcal{S} \otimes \mathcal{T}\), each of finite \(\mu \times \nu\) measure. By the definition of \(\sigma\)-finite (5.18), the measure \(\mu \times \nu\) is \(\sigma\)-finite.

Problem (5A.10)

Suppose \((X, \mathcal{S}, \mu)\) and \((Y, \mathcal{T}, \nu)\) are \(\sigma\)-finite measure spaces. Prove that if \(\omega\) is a measure on \(\mathcal{S} \otimes \mathcal{T}\) such that \(\omega(A \times B) = \mu(A)\nu(B)\) for all \(A \in \mathcal{S}\) and all \(B \in \mathcal{T}\), then \(\omega = \mu \times \nu\).

[The exercise above means that \(\mu \times \nu\) is the unique measure on \(\mathcal{S} \otimes \mathcal{T}\) that behaves as we expect on measurable rectangles.]

Solution

Throughout, \(\mu \times \nu\) is defined on \(\mathcal{S} \otimes \mathcal{T}\) by 5.25 (legitimate because both measure spaces are \(\sigma\)-finite) and is a measure by 5.27; thus the standard properties 2.57, 2.59, and 2.60 apply to \(\mu \times \nu\) as well as to \(\omega\).

Because \(\mu\) is \(\sigma\)-finite, Exercise 5A.8 provides an increasing sequence \(X_1 \subseteq X_2 \subseteq \cdots\) of sets in \(\mathcal{S}\) with \(X = \bigcup_{n=1}^{\infty} X_n\) and \(\mu(X_n) < \infty\) for every \(n \in \mathbf{Z}^+\). Similarly, there is an increasing sequence \(Y_1 \subseteq Y_2 \subseteq \cdots\) of sets in \(\mathcal{T}\) with \(Y = \bigcup_{n=1}^{\infty} Y_n\) and \(\nu(Y_n) < \infty\) for every \(n \in \mathbf{Z}^+\). For \(n \in \mathbf{Z}^+\) set

\begin{equation*} W_n = X_n \times Y_n . \end{equation*}

Then \(W_n\) is a measurable rectangle in \(\mathcal{S} \otimes \mathcal{T}\), and \(W_1 \subseteq W_2 \subseteq \cdots\) is an increasing sequence because both defining sequences are increasing. Moreover \(\bigcup_{n=1}^{\infty} W_n = X \times Y\): if \((x,y) \in X \times Y\), then \(x \in X_j\) and \(y \in Y_k\) for some \(j, k \in \mathbf{Z}^+\), and hence \((x,y) \in W_n\) with \(n = \max\{j, k\}\). Finally, by hypothesis and by 5.26,

\begin{equation*} \omega(W_n) = \mu(X_n)\, \nu(Y_n) = (\mu \times \nu)(W_n) < \infty \end{equation*}

for every \(n \in \mathbf{Z}^+\).

Fix \(n \in \mathbf{Z}^+\) and define

\begin{equation*} \mathcal{M}_n = \{ E \in \mathcal{S} \otimes \mathcal{T} : \omega(E \cap W_n) = (\mu \times \nu)(E \cap W_n) \} . \end{equation*}

We show that \(\mathcal{M}_n = \mathcal{S} \otimes \mathcal{T}\) by means of the Monotone Class Theorem (5.17).

Step 1: \(\mathcal{M}_n\) contains every measurable rectangle. Suppose \(A \in \mathcal{S}\) and \(B \in \mathcal{T}\). Then

\begin{equation*} (A \times B) \cap W_n = (A \times B) \cap (X_n \times Y_n) = (A \cap X_n) \times (B \cap Y_n), \end{equation*}

which is again a measurable rectangle in \(\mathcal{S} \otimes \mathcal{T}\). Hence, using the hypothesis on \(\omega\) and then 5.26,

\begin{equation*} \omega\bigl((A \times B) \cap W_n\bigr) = \mu(A \cap X_n)\, \nu(B \cap Y_n) = (\mu \times \nu)\bigl((A \times B) \cap W_n\bigr) . \end{equation*}

Thus \(A \times B \in \mathcal{M}_n\).

Step 2: \(\mathcal{M}_n\) contains the algebra \(\mathcal{A}\) of finite unions of measurable rectangles. Let \(\mathcal{A}\) denote the collection of finite unions of measurable rectangles in \(\mathcal{S} \otimes \mathcal{T}\); by 5.13(a) this is an algebra on \(X \times Y\). Suppose \(E \in \mathcal{A}\). By 5.13(b) we may write \(E = E_1 \cup \cdots \cup E_m\) where \(E_1, \ldots, E_m\) are disjoint measurable rectangles in \(\mathcal{S} \otimes \mathcal{T}\). The sets \(E_1 \cap W_n, \ldots, E_m \cap W_n\) are then disjoint sets in \(\mathcal{S} \otimes \mathcal{T}\) whose union is \(E \cap W_n\), so finite additivity of the measures \(\omega\) and \(\mu \times \nu\) (a measure is countably additive, and finite additivity follows by padding a finite disjoint list with copies of \(\varnothing\)) together with Step 1 gives

\begin{equation*} \omega(E \cap W_n) = \sum_{i=1}^{m} \omega(E_i \cap W_n) = \sum_{i=1}^{m} (\mu \times \nu)(E_i \cap W_n) = (\mu \times \nu)(E \cap W_n) . \end{equation*}

Thus \(E \in \mathcal{M}_n\), and hence \(\mathcal{A} \subseteq \mathcal{M}_n\).

Step 3: \(\mathcal{M}_n\) is a monotone class on \(X \times Y\). First suppose \(E_1 \subseteq E_2 \subseteq \cdots\) is an increasing sequence of sets in \(\mathcal{M}_n\), and let \(E = \bigcup_{k=1}^{\infty} E_k\), which is in \(\mathcal{S} \otimes \mathcal{T}\). Then \(E_1 \cap W_n \subseteq E_2 \cap W_n \subseteq \cdots\) is an increasing sequence with union \(E \cap W_n\), so 2.59, applied to \(\omega\) and then to \(\mu \times \nu\), gives

\begin{equation*} \omega(E \cap W_n) = \lim_{k \to \infty} \omega(E_k \cap W_n) = \lim_{k \to \infty} (\mu \times \nu)(E_k \cap W_n) = (\mu \times \nu)(E \cap W_n) . \end{equation*}

Hence \(E \in \mathcal{M}_n\).

Now suppose \(E_1 \supseteq E_2 \supseteq \cdots\) is a decreasing sequence of sets in \(\mathcal{M}_n\), and let \(E = \bigcap_{k=1}^{\infty} E_k \in \mathcal{S} \otimes \mathcal{T}\). Then \(E_1 \cap W_n \supseteq E_2 \cap W_n \supseteq \cdots\) is a decreasing sequence with intersection \(E \cap W_n\). Here is where the truncation to \(W_n\) matters: by 2.57(a),

\begin{equation*} \omega(E_1 \cap W_n) \le \omega(W_n) < \infty \quad\text{and}\quad (\mu \times \nu)(E_1 \cap W_n) \le (\mu \times \nu)(W_n) < \infty , \end{equation*}

so 2.60 applies to both measures and yields

\begin{equation*} \omega(E \cap W_n) = \lim_{k \to \infty} \omega(E_k \cap W_n) = \lim_{k \to \infty} (\mu \times \nu)(E_k \cap W_n) = (\mu \times \nu)(E \cap W_n) . \end{equation*}

Hence \(E \in \mathcal{M}_n\). Thus \(\mathcal{M}_n\) is a monotone class.

Step 4: conclusion for fixed \(n\). By Steps 2 and 3, \(\mathcal{M}_n\) is a monotone class containing the algebra \(\mathcal{A}\), so \(\mathcal{M}_n\) contains the smallest monotone class containing \(\mathcal{A}\), which by the Monotone Class Theorem (5.17) equals the smallest \(\sigma\)-algebra containing \(\mathcal{A}\). That smallest \(\sigma\)-algebra is \(\mathcal{S} \otimes \mathcal{T}\): indeed \(\mathcal{A}\) contains all measurable rectangles, so any \(\sigma\)-algebra containing \(\mathcal{A}\) contains \(\mathcal{S} \otimes \mathcal{T}\) by 5.2, while conversely \(\mathcal{A} \subseteq \mathcal{S} \otimes \mathcal{T}\) since \(\mathcal{S} \otimes \mathcal{T}\) is a \(\sigma\)-algebra containing the measurable rectangles. Because also \(\mathcal{M}_n \subseteq \mathcal{S} \otimes \mathcal{T}\) by definition, we conclude

\begin{equation*} \mathcal{M}_n = \mathcal{S} \otimes \mathcal{T} . \end{equation*}

In other words, \(\omega(E \cap W_n) = (\mu \times \nu)(E \cap W_n)\) for every \(E \in \mathcal{S} \otimes \mathcal{T}\) and every \(n \in \mathbf{Z}^+\).

Step 5: removing the truncation. Suppose \(E \in \mathcal{S} \otimes \mathcal{T}\). Then \(E \cap W_1 \subseteq E \cap W_2 \subseteq \cdots\) is an increasing sequence of sets in \(\mathcal{S} \otimes \mathcal{T}\) whose union is

\begin{equation*} \bigcup_{n=1}^{\infty} (E \cap W_n) = E \cap \Bigl( \bigcup_{n=1}^{\infty} W_n \Bigr) = E \cap (X \times Y) = E . \end{equation*}

Applying 2.59 to \(\omega\) and to \(\mu \times \nu\), and using Step 4,

\begin{equation*} \omega(E) = \lim_{n \to \infty} \omega(E \cap W_n) = \lim_{n \to \infty} (\mu \times \nu)(E \cap W_n) = (\mu \times \nu)(E) . \end{equation*}

Because \(E \in \mathcal{S} \otimes \mathcal{T}\) was arbitrary, \(\omega = \mu \times \nu\), as desired.

Exercises 5B

Problem (5B.1)

(a) Let \(\lambda\) denote Lebesgue measure on \([0,1]\). Show that

\begin{equation*} \int_{[0,1]}\int_{[0,1]}\frac{x^2-y^2}{(x^2+y^2)^2}\,d\lambda(y)\,d\lambda(x)=\frac{\pi}{4} \end{equation*}

and

\begin{equation*} \int_{[0,1]}\int_{[0,1]}\frac{x^2-y^2}{(x^2+y^2)^2}\,d\lambda(x)\,d\lambda(y)=-\frac{\pi}{4}. \end{equation*}

(b) Explain why (a) violates neither Tonelli’s Theorem nor Fubini’s Theorem.

Solution

Throughout, let

\begin{equation*} f(x,y)=\frac{x^2-y^2}{(x^2+y^2)^2}\quad\text{for }(x,y)\in[0,1]^2\setminus\{(0,0)\}, \qquad f(0,0)=0 . \end{equation*}

(The value assigned at the origin is irrelevant for everything below.) Note the antisymmetry \(f(y,x)=-f(x,y)\).

(a) The key computation is the pair of antiderivatives

\begin{equation*} \frac{\partial}{\partial y}\Bigl(\frac{y}{x^2+y^2}\Bigr) =\frac{(x^2+y^2)-y\cdot 2y}{(x^2+y^2)^2}=\frac{x^2-y^2}{(x^2+y^2)^2}, \qquad \frac{\partial}{\partial x}\Bigl(\frac{-x}{x^2+y^2}\Bigr) =\frac{x^2-y^2}{(x^2+y^2)^2}, \end{equation*}

valid at every point where \((x,y)\neq(0,0)\).

Fix \(x\in(0,1]\). Then \(y\mapsto f(x,y)\) is continuous on \([0,1]\) (the denominator is at least \(x^4>0\) there), hence Riemann integrable, and by 3.34 its Lebesgue integral equals its Riemann integral. The Fundamental Theorem of Calculus and the first antiderivative above give

\begin{equation*} \int_{[0,1]}f(x,y)\,d\lambda(y)=\left[\frac{y}{x^2+y^2}\right]_{y=0}^{y=1}=\frac{1}{1+x^2}. \end{equation*}

For \(x=0\) we have \(f(0,y)=-1/y^2\) for \(y\in(0,1]\) and \(f(0,0)=0\); thus \(f(0,\cdot)\) is Borel measurable, with \(f(0,\cdot)^+=0\) and \(f(0,\cdot)^-=y^{-2}\chi_{(0,1]}\). For each \(k\in\mathbb{Z}^+\) the integrand \(y^{-2}\) is continuous on \([1/k,1]\), so 3.34 gives

\begin{equation*} \int_{[0,1]}f(0,\cdot)^-\,d\lambda\ \ge\ \int_{[1/k,1]}\frac{dy}{y^2}=k-1 , \end{equation*}

and letting \(k\to\infty\) shows \(\int_{[0,1]}f(0,\cdot)^-\,d\lambda=\infty\). Since \(\int_{[0,1]}f(0,\cdot)^+\,d\lambda=0\) is finite, the integral is defined by 3.18 and equals \(0-\infty=-\infty\).

Thus the inner integral, as a function of \(x\), is the function \(g:[0,1]\to[-\infty,\infty]\) given by \(g(0)=-\infty\) and \(g(x)=1/(1+x^2)\) for \(x>0\). This \(g\) is \(\mathcal{B}\)-measurable, since for \(a\in\mathbb{R}\) the set \(g^{-1}((a,\infty])\) equals \((0,1]\) intersected with an interval (the map \(x\mapsto 1/(1+x^2)\) is decreasing on \((0,1]\), and \(0\) never belongs to this set because \(g(0)=-\infty\)); so 2.52 applies. Moreover \(g^-=\infty\cdot\chi_{\{0\}}\), and \(\int_{[0,1]}g^-\,d\lambda=\infty\cdot\lambda(\{0\})=\infty\cdot 0=0\), so 3.18 gives \(\int_{[0,1]}g\,d\lambda=\int_{[0,1]}g^+\,d\lambda\). Now \(g^+\) differs from the continuous function \(x\mapsto 1/(1+x^2)\) only at the point \(0\); since the integral of a nonnegative measurable function over a set of measure \(0\) is \(0\), additivity of integration over the decomposition \([0,1]=\{0\}\cup(0,1]\) shows that \(g^+\) and that continuous function have the same integral over \([0,1]\). Hence, using 3.34 again,

\begin{equation*} \int_{[0,1]}g\,d\lambda=\int_{[0,1]}g^+\,d\lambda=\int_0^1\frac{dx}{1+x^2}=\arctan 1-\arctan 0=\frac{\pi}{4}. \end{equation*}

This proves the first equation:

\begin{equation*} \int_{[0,1]}\int_{[0,1]}\frac{x^2-y^2}{(x^2+y^2)^2}\,d\lambda(y)\,d\lambda(x)=\frac{\pi}{4}. \end{equation*}

For the second equation, fix \(y\in(0,1]\). Again \(x\mapsto f(x,y)\) is continuous on \([0,1]\), and the second antiderivative above gives

\begin{equation*} \int_{[0,1]}f(x,y)\,d\lambda(x)=\left[\frac{-x}{x^2+y^2}\right]_{x=0}^{x=1}=-\frac{1}{1+y^2}. \end{equation*}

For \(y=0\) we get \(f(x,0)=1/x^2\) for \(x\in(0,1]\), so by the same estimate as before this inner integral equals \(+\infty\). Hence the inner integral, as a function of \(y\), is \(h:[0,1]\to[-\infty,\infty]\) with \(h(0)=\infty\) and \(h(y)=-1/(1+y^2)\) for \(y>0\); as above \(h\) is \(\mathcal{B}\)-measurable, \(h^+=\infty\cdot\chi_{\{0\}}\) has integral \(\infty\cdot 0=0\), and \(h^-\) differs only at the point \(0\) from the continuous function \(y\mapsto 1/(1+y^2)\). So 3.18 and 3.34 give

\begin{equation*} \int_{[0,1]}\int_{[0,1]}\frac{x^2-y^2}{(x^2+y^2)^2}\,d\lambda(x)\,d\lambda(y)=0-\int_0^1\frac{dy}{1+y^2}=-\frac{\pi}{4}. \end{equation*}

(Alternatively this follows from the first equation and the antisymmetry \(f(y,x)=-f(x,y)\), which interchanges the roles of the two variables.)

(b) First note that \(f\) is \(\mathcal{B}\otimes\mathcal{B}\)-measurable on \([0,1]^2\), where \(\mathcal{B}\) denotes the Borel subsets of \([0,1]\). Indeed, the coordinate functions \(\pi_1(x,y)=x\) and \(\pi_2(x,y)=y\) are \(\mathcal{B}\otimes\mathcal{B}\)-measurable, because for each \(a\in\mathbb{R}\) the sets

\begin{equation*} \pi_1^{-1}\bigl((a,\infty)\bigr)=\bigl((a,\infty)\cap[0,1]\bigr)\times[0,1], \qquad \pi_2^{-1}\bigl((a,\infty)\bigr)=[0,1]\times\bigl((a,\infty)\cap[0,1]\bigr) \end{equation*}

are measurable rectangles, so 2.39 applies. Hence by 2.46 the functions \((x,y)\mapsto x^2-y^2\) and \((x,y)\mapsto(x^2+y^2)^2\) are \(\mathcal{B}\otimes\mathcal{B}\)-measurable. The denominator vanishes at the origin, so to stay within 2.46(b) define, for \(k\in\mathbb{Z}^+\),

\begin{equation*} f_k(x,y)=\frac{x^2-y^2}{(x^2+y^2)^2+\frac{1}{k}} . \end{equation*}

The denominator here is at least \(1/k>0\) everywhere, so each \(f_k\) is a real-valued \(\mathcal{B}\otimes\mathcal{B}\)-measurable function by 2.46. If \((x,y)\neq(0,0)\) then \((x^2+y^2)^2>0\) and \(f_k(x,y)\to f(x,y)\); and \(f_k(0,0)=0=f(0,0)\) for every \(k\). Thus \(f_k\to f\) pointwise on all of \([0,1]^2\), so \(f\) is \(\mathcal{B}\otimes\mathcal{B}\)-measurable by 2.48.

Tonelli’s Theorem (5.28) is not contradicted, because it requires the function being integrated to be nonnegative, and \(f\) takes both positive values (where \(|x|>|y|\)) and negative values (where \(|x|<|y|\)). The other hypotheses of 5.28 do hold here: \(\lambda\) restricted to \([0,1]\) is a finite measure, hence \(\sigma\)-finite. So this exercise shows that the nonnegativity hypothesis in Tonelli’s Theorem cannot be dropped.

Fubini’s Theorem (5.32) is not contradicted, because its hypothesis \(\int_{[0,1]^2}|f|\,d(\lambda\times\lambda)<\infty\) fails. To see this directly, apply Tonelli’s Theorem to the nonnegative measurable function \(|f|\):

\begin{equation*} \int_{[0,1]^2}|f|\,d(\lambda\times\lambda)=\int_{[0,1]}\int_{[0,1]}|f(x,y)|\,d\lambda(y)\,d\lambda(x). \end{equation*}

Fix \(x\in(0,1]\). On \(\{y:0\le y\le x\}\) we have \(f(x,y)\ge 0\), so

\begin{equation*} \int_{[0,1]}|f(x,y)|\,d\lambda(y)\ge\int_{[0,x]}f(x,y)\,d\lambda(y) =\left[\frac{y}{x^2+y^2}\right]_{y=0}^{y=x}=\frac{x}{2x^2}=\frac{1}{2x}. \end{equation*}

Therefore

\begin{equation*} \int_{[0,1]^2}|f|\,d(\lambda\times\lambda)\ge\int_{(0,1]}\frac{1}{2x}\,d\lambda(x) \ge\frac{1}{2}\int_{[1/k,\,1]}\frac{dx}{x}=\frac{\ln k}{2} \end{equation*}

for every \(k\in\mathbb{Z}^+\) (the middle integrals being computed by 3.34, the integrand being continuous on \([1/k,1]\)), and letting \(k\to\infty\) gives \(\int_{[0,1]^2}|f|\,d(\lambda\times\lambda)=\infty\). So \(f\notin\mathcal{L}^1(\lambda\times\lambda)\) and Fubini’s Theorem does not apply.

(Conversely, once one knows the two iterated integrals in (a) are unequal, Fubini’s Theorem itself forces \(\int|f|\,d(\lambda\times\lambda)=\infty\); the computation above verifies this independently.)

Problem (5B.2)

(a) Give an example of a doubly indexed collection \(\{x_{m,n}:m,n\in\mathbb{Z}^+\}\) of real numbers such that

\begin{equation*} \sum_{m=1}^{\infty}\sum_{n=1}^{\infty}x_{m,n}=0 \qquad\text{and}\qquad \sum_{n=1}^{\infty}\sum_{m=1}^{\infty}x_{m,n}=\infty. \end{equation*}

(b) Explain why (a) violates neither Tonelli’s Theorem nor Fubini’s Theorem.

Solution

(a) Define

\begin{equation*} x_{m,n}= \begin{cases} 1 & \text{if }n=m,\\ -1 & \text{if }n=2m,\\ 0 & \text{otherwise.} \end{cases} \end{equation*}

Thus each row of the array consists of a \(+1\) on the diagonal and a \(-1\) placed twice as far out; the \(-1\) of row \(m\) lands in an even column, so the rows cancel but the odd columns never receive a compensating \(-1\).

Row sums. Fix \(m\in\mathbb{Z}^+\). Since \(m\ge 1\), we have \(2m\neq m\), so the sequence \(n\mapsto x_{m,n}\) has exactly two nonzero terms, namely \(x_{m,m}=1\) and \(x_{m,2m}=-1\). Hence

\begin{equation*} \sum_{n=1}^{\infty}x_{m,n}=1+(-1)=0 \qquad\text{for every }m\in\mathbb{Z}^+, \end{equation*}

and therefore

\begin{equation*} \sum_{m=1}^{\infty}\sum_{n=1}^{\infty}x_{m,n}=\sum_{m=1}^{\infty}0=0. \end{equation*}

Column sums. Fix \(n\in\mathbb{Z}^+\). The term \(x_{m,n}\) is nonzero only if \(m=n\) (giving \(+1\)) or \(2m=n\) (giving \(-1\), which requires \(n\) to be even, with \(m=n/2\)); these two indices are distinct because \(n\neq n/2\) for \(n\ge 1\). Hence

\begin{equation*} \sum_{m=1}^{\infty}x_{m,n}= \begin{cases} 1 & \text{if \(n\) is odd},\\ 1+(-1)=0 & \text{if \(n\) is even.} \end{cases} \end{equation*}

The partial sums of \(\sum_{n}\bigl(\sum_{m}x_{m,n}\bigr)\) therefore count the odd integers up to \(n\) and increase without bound, so

\begin{equation*} \sum_{n=1}^{\infty}\sum_{m=1}^{\infty}x_{m,n}=\infty . \end{equation*}

(b) Put this in the measure-theoretic setting. Let \(X=Y=\mathbb{Z}^+\), let \(\mathcal{S}=\mathcal{T}\) be the \(\sigma\)-algebra of all subsets of \(\mathbb{Z}^+\), and let \(\mu=\nu\) be counting measure. These measure spaces are \(\sigma\)-finite, since \(\mathbb{Z}^+=\bigcup_{k=1}^{\infty}\{1,\dots,k\}\) and each \(\{1,\dots,k\}\) has finite measure. Define \(f:\mathbb{Z}^+\times\mathbb{Z}^+\to\mathbb{R}\) by \(f(m,n)=x_{m,n}\). Every subset of \(\mathbb{Z}^+\times\mathbb{Z}^+\) belongs to \(\mathcal{S}\otimes\mathcal{T}\), because each singleton \(\{(m,n)\}=\{m\}\times\{n\}\) is a measurable rectangle and every subset of the countable set \(\mathbb{Z}^+\times\mathbb{Z}^+\) is a countable union of singletons; hence \(f\) is \(\mathcal{S}\otimes\mathcal{T}\)-measurable. Integration with respect to counting measure is summation (3.6, and 3.42 for the real-valued case), so the two iterated sums in (a) are exactly the two iterated integrals appearing in 5.28 and 5.32. All the inner integrals here are defined: for fixed \(m\) the function \(n\mapsto x_{m,n}\) has only two nonzero values, and for fixed \(n\) the function \(m\mapsto x_{m,n}\) has at most two nonzero values, so in each case the function lies in \(\mathcal{L}^1\) of counting measure.

Tonelli’s Theorem (5.28) is not contradicted, because it requires the integrand to be nonnegative; our \(f\) takes the value \(-1\). The same remark applies to its corollary 5.31 on double sums, whose hypothesis is precisely that the numbers \(x_{m,n}\) be nonnegative. Note that the \(\sigma\)-finiteness hypothesis does hold here, so nonnegativity is the only hypothesis that fails; the example thus shows that nonnegativity cannot be dropped from 5.31.

Fubini’s Theorem (5.32) is not contradicted, because its hypothesis \(\int_{\mathbb{Z}^+\times\mathbb{Z}^+}|f|\,d(\mu\times\nu)<\infty\) fails. Indeed, \(|f|\) is nonnegative and measurable, so Tonelli’s Theorem applies to it and gives

\begin{equation*} \int_{\mathbb{Z}^+\times\mathbb{Z}^+}|f|\,d(\mu\times\nu) =\sum_{m=1}^{\infty}\sum_{n=1}^{\infty}|x_{m,n}| =\sum_{m=1}^{\infty}2=\infty , \end{equation*}

since each row has exactly two nonzero entries, each of absolute value \(1\). So \(f\) is not integrable with respect to \(\mu\times\nu\) and Fubini’s Theorem says nothing about it.

Problem (5B.3)

Suppose \((X,\mathcal{S})\) is a measurable space and \(f:X\to[0,\infty]\) is a function. Let \(\mathcal{B}\) denote the \(\sigma\)-algebra of Borel subsets of \((0,\infty)\). Prove that \(U_f\in\mathcal{S}\otimes\mathcal{B}\) if and only if \(f\) is an \(\mathcal{S}\)-measurable function.

Solution

Recall from 5.34 that

\begin{equation*} U_f=\{(x,t)\in X\times(0,\infty):0<t<f(x)\}. \end{equation*}

Note that no measure is involved in this exercise; only the \(\sigma\)-algebras matter.

First suppose \(f\) is an \(\mathcal{S}\)-measurable function. We claim that

\begin{equation*} U_f=\bigcup_{r\in\mathbb{Q}\cap(0,\infty)}f^{-1}\bigl((r,\infty]\bigr)\times(0,r]. \end{equation*}

To prove the inclusion \(\supseteq\), suppose \(r\in\mathbb{Q}\cap(0,\infty)\), \(x\in f^{-1}((r,\infty])\), and \(t\in(0,r]\). Then \(0<t\le r<f(x)\), so \((x,t)\in U_f\). To prove the inclusion \(\subseteq\), suppose \((x,t)\in U_f\), so \(0<t<f(x)\). Because \(\mathbb{Q}\) is dense in \(\mathbb{R}\) there exists \(r\in\mathbb{Q}\) with \(t<r<f(x)\) (if \(f(x)=\infty\), take any rational \(r>t\)). Then \(r>t>0\), so \(r\in\mathbb{Q}\cap(0,\infty)\), and \(x\in f^{-1}((r,\infty])\) and \(t\in(0,r]\). This proves the claim.

Now for each \(r\in\mathbb{Q}\cap(0,\infty)\) the set \((r,\infty]\) is a Borel subset of \([-\infty,\infty]\) (see 2.50), so \(f^{-1}((r,\infty])\in\mathcal{S}\) because \(f\) is \(\mathcal{S}\)-measurable (2.51); also \((0,r]\) is a Borel subset of \((0,\infty)\), so \((0,r]\in\mathcal{B}\). Hence each set \(f^{-1}((r,\infty])\times(0,r]\) is a measurable rectangle in \(\mathcal{S}\otimes\mathcal{B}\). The displayed union is over the countable set \(\mathbb{Q}\cap(0,\infty)\), and \(\mathcal{S}\otimes\mathcal{B}\) is closed under countable unions. Thus \(U_f\in\mathcal{S}\otimes\mathcal{B}\). (This is the same conclusion reached, by a slightly different decomposition, in the first paragraph of the proof of 5.35; the argument there also uses no properties of the measure.)

Conversely, suppose \(U_f\in\mathcal{S}\otimes\mathcal{B}\). Fix \(t\in(0,\infty)\). Because cross sections of measurable sets are measurable (5.6), the cross section

\begin{equation*} [U_f]^t=\{x\in X:(x,t)\in U_f\} \end{equation*}

belongs to \(\mathcal{S}\). Since \(t>0\), the condition \((x,t)\in U_f\) says exactly that \(t<f(x)\). In other words,

\begin{equation*} f^{-1}\bigl((t,\infty]\bigr)=\{x\in X:f(x)>t\}=[U_f]^t\in\mathcal{S} \qquad\text{for every }t\in(0,\infty). \end{equation*}

Now let \(a\in\mathbb{R}\) be arbitrary; we check that \(f^{-1}((a,\infty])\in\mathcal{S}\) in all three possible cases.

  • If \(a>0\), then \(f^{-1}((a,\infty])=[U_f]^a\in\mathcal{S}\) by the previous paragraph.

  • If \(a<0\), then \(f^{-1}((a,\infty])=X\in\mathcal{S}\), because \(f\) takes values in \([0,\infty]\).

  • If \(a=0\), then \(f(x)>0\) if and only if \(f(x)>1/k\) for some \(k\in\mathbb{Z}^+\), so

\begin{equation*} f^{-1}\bigl((0,\infty]\bigr)=\bigcup_{k=1}^{\infty}f^{-1}\bigl((\tfrac{1}{k},\infty]\bigr) =\bigcup_{k=1}^{\infty}[U_f]^{1/k}, \end{equation*}

which is a countable union of elements of \(\mathcal{S}\) and hence belongs to \(\mathcal{S}\).

Thus \(f^{-1}((a,\infty])\in\mathcal{S}\) for every \(a\in\mathbb{R}\), and therefore \(f\) is an \(\mathcal{S}\)-measurable function by 2.52.

Problem (5B.4)

Suppose \((X,\mathcal{S})\) is a measurable space and \(f:X\to\mathbb{R}\) is a function. Let \(\operatorname{graph}(f)\subseteq X\times\mathbb{R}\) denote the graph of \(f\):

\begin{equation*} \operatorname{graph}(f)=\bigl\{\bigl(x,f(x)\bigr):x\in X\bigr\}. \end{equation*}

Let \(\mathcal{B}\) denote the \(\sigma\)-algebra of Borel subsets of \(\mathbb{R}\). Prove that \(\operatorname{graph}(f)\in\mathcal{S}\otimes\mathcal{B}\) if \(f\) is an \(\mathcal{S}\)-measurable function.

Solution

Suppose \(f\) is an \(\mathcal{S}\)-measurable function.

Define \(g,h:X\times\mathbb{R}\to\mathbb{R}\) by

\begin{equation*} g(x,y)=f(x)\qquad\text{and}\qquad h(x,y)=y . \end{equation*}

Both are \(\mathcal{S}\otimes\mathcal{B}\)-measurable. Indeed, if \(B\subseteq\mathbb{R}\) is a Borel set, then

\begin{equation*} g^{-1}(B)=f^{-1}(B)\times\mathbb{R} \qquad\text{and}\qquad h^{-1}(B)=X\times B , \end{equation*}

and \(f^{-1}(B)\in\mathcal{S}\) because \(f\) is \(\mathcal{S}\)-measurable, so both of these sets are measurable rectangles in \(\mathcal{S}\otimes\mathcal{B}\).

By 2.46(a), the function \(g-h\) is \(\mathcal{S}\otimes\mathcal{B}\)-measurable, where

\begin{equation*} (g-h)(x,y)=f(x)-y . \end{equation*}

Because \(\{0\}\) is a Borel subset of \(\mathbb{R}\), we conclude that

\begin{equation*} \operatorname{graph}(f)=\{(x,y)\in X\times\mathbb{R}:f(x)-y=0\}=(g-h)^{-1}\bigl(\{0\}\bigr)\in\mathcal{S}\otimes\mathcal{B}, \end{equation*}

as desired.

For a proof that avoids citing 2.46 and instead exhibits \(\operatorname{graph}(f)\) explicitly from measurable rectangles, note that

\begin{equation*} \operatorname{graph}(f)=\bigcap_{k=1}^{\infty}\ \bigcup_{j\in\mathbb{Z}} \Bigl(f^{-1}\bigl([\tfrac{j}{k},\tfrac{j+1}{k})\bigr)\times[\tfrac{j}{k},\tfrac{j+1}{k})\Bigr). \end{equation*}

To verify the inclusion \(\subseteq\), suppose \(y=f(x)\) and let \(k\in\mathbb{Z}^+\). Take \(j\) to be the unique integer with \(j\le kf(x)<j+1\); then \(f(x)\in[\tfrac{j}{k},\tfrac{j+1}{k})\), so \(x\in f^{-1}([\tfrac{j}{k},\tfrac{j+1}{k}))\) and also \(y=f(x)\in[\tfrac{j}{k},\tfrac{j+1}{k})\). Hence \((x,y)\) lies in the \(k^{\text{th}}\) union, for every \(k\).

To verify the inclusion \(\supseteq\), suppose \((x,y)\) lies in the set on the right. For each \(k\in\mathbb{Z}^+\) there exists \(j_k\in\mathbb{Z}\) with \(f(x)\in[\tfrac{j_k}{k},\tfrac{j_k+1}{k})\) and \(y\in[\tfrac{j_k}{k},\tfrac{j_k+1}{k})\). Both numbers lie in an interval of length \(\tfrac{1}{k}\), so \(|y-f(x)|<\tfrac{1}{k}\). As this holds for every \(k\in\mathbb{Z}^+\), we get \(y=f(x)\), so \((x,y)\in\operatorname{graph}(f)\).

Finally, each interval \([\tfrac{j}{k},\tfrac{j+1}{k})\) is a Borel subset of \(\mathbb{R}\), so \(f^{-1}([\tfrac{j}{k},\tfrac{j+1}{k}))\in\mathcal{S}\) and each set \(f^{-1}([\tfrac{j}{k},\tfrac{j+1}{k}))\times[\tfrac{j}{k},\tfrac{j+1}{k})\) is a measurable rectangle in \(\mathcal{S}\otimes\mathcal{B}\). The union over \(j\in\mathbb{Z}\) is a countable union and the intersection over \(k\in\mathbb{Z}^+\) is a countable intersection, and a \(\sigma\)-algebra is closed under both operations. Hence \(\operatorname{graph}(f)\in\mathcal{S}\otimes\mathcal{B}\).

Exercises 5C

Problem (5C.1)

Show that a set \(G \subseteq \mathbb{R}^n\) is open in \(\mathbb{R}^n\) if and only if for each \((b_1, \ldots, b_n) \in G\), there exists \(r > 0\) such that

\begin{equation*} \{ (a_1, \ldots, a_n) \in \mathbb{R}^n : \sqrt{(a_1 - b_1)^2 + \cdots + (a_n - b_n)^2} < r \} \subseteq G. \end{equation*}

Solution

Recall the definitions used in this section: \(\|(x_1,\ldots,x_n)\|_\infty = \max\{|x_1|,\ldots,|x_n|\}\), the open cube \(B(x,\delta) = \{ y \in \mathbb{R}^n : \|y - x\|_\infty < \delta \}\), and \(G \subseteq \mathbb{R}^n\) is open means that for every \(x \in G\) there exists \(\delta > 0\) with \(B(x, \delta) \subseteq G\).

Write \(\|x\|_2 = \sqrt{x_1^2 + \cdots + x_n^2}\) for the Euclidean norm, and for \(b \in \mathbb{R}^n\) and \(r>0\) let

\begin{equation*} D(b, r) = \{ a \in \mathbb{R}^n : \|a - b\|_2 < r \} \end{equation*}

denote the Euclidean open ball; this is the set displayed in the statement of the exercise. So we must show that \(G\) is open if and only if every \(b \in G\) has some \(r > 0\) with \(D(b,r) \subseteq G\).

The proof rests on the elementary comparison of the two norms:

\begin{equation*} \|x\|_\infty \le \|x\|_2 \le \sqrt{n}\, \|x\|_\infty \quad \text{for all } x \in \mathbb{R}^n. \end{equation*}

To verify the first inequality, choose \(j\) with \(|x_j| = \|x\|_\infty\); then \(\|x\|_\infty^2 = x_j^2 \le x_1^2 + \cdots + x_n^2 = \|x\|_2^2\), and taking square roots gives \(\|x\|_\infty \le \|x\|_2\). For the second inequality, \(x_k^2 \le \|x\|_\infty^2\) for each \(k\), so \(\|x\|_2^2 = x_1^2 + \cdots + x_n^2 \le n \|x\|_\infty^2\), and taking square roots gives \(\|x\|_2 \le \sqrt{n}\, \|x\|_\infty\).

These two inequalities say exactly that

\begin{equation*} D(b, \delta) \subseteq B(b, \delta) \quad \text{and} \quad B(b, r/\sqrt{n}\,) \subseteq D(b, r) \end{equation*}

for all \(b \in \mathbb{R}^n\), \(\delta > 0\), and \(r > 0\). Indeed, if \(\|a - b\|_2 < \delta\) then \(\|a-b\|_\infty \le \|a-b\|_2 < \delta\), which gives the first inclusion; and if \(\|a - b\|_\infty < r/\sqrt{n}\) then \(\|a - b\|_2 \le \sqrt{n}\,\|a-b\|_\infty < r\), which gives the second.

Now suppose \(G\) is open and \(b = (b_1,\ldots,b_n) \in G\). Then there exists \(\delta > 0\) with \(B(b, \delta) \subseteq G\). Taking \(r = \delta\), the first inclusion above gives

\begin{equation*} D(b, r) = D(b,\delta) \subseteq B(b, \delta) \subseteq G, \end{equation*}

which is the desired conclusion.

Conversely, suppose that for each \(b \in G\) there exists \(r > 0\) with \(D(b, r) \subseteq G\). Let \(b \in G\) and choose such an \(r\). Put \(\delta = r/\sqrt{n} > 0\). Then the second inclusion above gives

\begin{equation*} B(b, \delta) = B(b, r/\sqrt{n}\,) \subseteq D(b, r) \subseteq G. \end{equation*}

Thus every point of \(G\) is the center of an open cube contained in \(G\), which means that \(G\) is open.

Hence the two conditions are equivalent. (In other words, the topology on \(\mathbb{R}^n\) generated by cubes, which is the one used throughout this section, is the same as the usual topology generated by Euclidean balls; equivalently, \(\|\cdot\|_\infty\) and \(\|\cdot\|_2\) are equivalent norms on \(\mathbb{R}^n\).)

Problem (5C.2)

Show that there exists a set \(E \subseteq \mathbb{R}^2\) (thinking of \(\mathbb{R}^2\) as equal to \(\mathbb{R} \times \mathbb{R}\)) such that the cross sections \([E]_a\) and \([E]^a\) are open subsets of \(\mathbb{R}\) for every \(a \in \mathbb{R}\), but \(E \notin \mathcal{B}_2\).

Solution

Recall from 5.4 that \([E]_a = \{ y \in \mathbb{R} : (a,y) \in E \}\) and \([E]^a = \{ x \in \mathbb{R} : (x,a) \in E \}\).

The strategy is a counting argument: we exhibit a family of \(2^{\mathfrak{c}}\) distinct subsets of \(\mathbb{R}^2\) all of whose cross sections are open (here \(\mathfrak{c} = 2^{\aleph_0}\) is the cardinality of \(\mathbb{R}\)), and we show that \(\mathcal{B}_2\) has only \(\mathfrak{c}\) elements. Since \(2^{\mathfrak{c}} > \mathfrak{c}\) by Cantor’s Theorem, some member of the family is not a Borel set.

Step 1: a large family of sets with open cross sections.

For each \(A \subseteq (0, \infty)\) define \(f_A : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} f_A(x) = \begin{cases} -x & \text{if } |x| \in A, \\ x & \text{if } |x| \notin A. \end{cases} \end{equation*}

Because \(|f_A(x)| = |x|\) for every \(x \in \mathbb{R}\), we have \(f_A\bigl(f_A(x)\bigr) = x\) for every \(x\); thus \(f_A\) is a bijection of \(\mathbb{R}\) onto \(\mathbb{R}\) that is its own inverse. Let

\begin{equation*} \Gamma_A = \{ (x, f_A(x)) : x \in \mathbb{R} \} \quad \text{and} \quad E_A = \mathbb{R}^2 \setminus \Gamma_A . \end{equation*}

Fix \(a \in \mathbb{R}\). Then \((a, y) \in \Gamma_A\) if and only if \(y = f_A(a)\), so

\begin{equation*} [E_A]_a = \mathbb{R} \setminus \{ f_A(a) \}, \end{equation*}

which is an open subset of \(\mathbb{R}\). Also \((x, a) \in \Gamma_A\) if and only if \(f_A(x) = a\), which (because \(f_A\) is a bijection equal to its own inverse) happens if and only if \(x = f_A(a)\). Hence

\begin{equation*} [E_A]^a = \mathbb{R} \setminus \{ f_A(a) \}, \end{equation*}

which is again an open subset of \(\mathbb{R}\). So every set \(E_A\) has all its cross sections open.

The map \(A \mapsto E_A\) is injective. Indeed, \(E_A\) determines \(\Gamma_A = \mathbb{R}^2 \setminus E_A\), and for \(t > 0\) we have \(f_A(t) = -t\) if \(t \in A\) and \(f_A(t) = t\) if \(t \notin A\); because \(-t \ne t\) for \(t > 0\), this gives

\begin{equation*} A = \{ t > 0 : (t, -t) \in \Gamma_A \}, \end{equation*}

so \(A\) is recovered from \(E_A\). Since \(|(0,\infty)| = \mathfrak{c}\), the collection \(\{ E_A : A \subseteq (0,\infty) \}\) consists of \(2^{\mathfrak{c}}\) distinct subsets of \(\mathbb{R}^2\), each of which has all cross sections open.

Step 2: \(\mathcal{B}_2\) has cardinality at most \(\mathfrak{c}\).

Let \(\Sigma_0\) denote the collection of open subsets of \(\mathbb{R}^2\). By the proof of 5.38(a), every open subset of \(\mathbb{R}^2\) is the union of a subcollection of the countable collection \(\mathcal{Q}\) of open cubes whose center has rational coordinates and whose side length is rational. The map sending a subcollection of \(\mathcal{Q}\) to its union therefore maps the power set of \(\mathcal{Q}\) onto \(\Sigma_0\), so \(|\Sigma_0| \le 2^{\aleph_0} = \mathfrak{c}\).

Let \(\omega_1\) denote the first uncountable ordinal. For each ordinal \(\alpha\) with \(0 < \alpha < \omega_1\) define

\begin{equation*} \Sigma_\alpha = \Bigl\{ \bigcup_{k=1}^{\infty} F_k : \text{for each } k, \ F_k \in \textstyle\bigcup_{\beta < \alpha} \Sigma_\beta \text{ or } \mathbb{R}^2 \setminus F_k \in \textstyle\bigcup_{\beta < \alpha} \Sigma_\beta \Bigr\}, \end{equation*}

and let \(\Sigma = \bigcup_{\alpha < \omega_1} \Sigma_\alpha\). Taking \(F_1 = F_2 = \cdots = E\) shows that \(\Sigma_\beta \subseteq \Sigma_\alpha\) whenever \(\beta < \alpha < \omega_1\).

The collection \(\Sigma\) is a \(\sigma\)-algebra on \(\mathbb{R}^2\). It is closed under complementation: if \(E \in \Sigma_\alpha\), then taking \(F_1 = F_2 = \cdots = \mathbb{R}^2 \setminus E\) (whose complement \(E\) lies in \(\Sigma_\alpha\)) shows \(\mathbb{R}^2 \setminus E \in \Sigma_{\alpha + 1} \subseteq \Sigma\). It is closed under countable unions: if \(E_1, E_2, \ldots \in \Sigma\), choose \(\alpha_k < \omega_1\) with \(E_k \in \Sigma_{\alpha_k}\); the supremum \(\alpha\) of the countable set \(\{\alpha_k\}\) is a countable ordinal, so \(\alpha < \omega_1\), and then every \(E_k\) lies in \(\Sigma_\alpha\), whence \(\bigcup_k E_k \in \Sigma_{\alpha+1} \subseteq \Sigma\). Finally \(\Sigma_0 \subseteq \Sigma\), so \(\Sigma\) is a \(\sigma\)-algebra containing every open subset of \(\mathbb{R}^2\); therefore \(\mathcal{B}_2 \subseteq \Sigma\).

Now we bound \(|\Sigma|\) by transfinite induction. We have \(|\Sigma_0| \le \mathfrak{c}\). Suppose \(0 < \alpha < \omega_1\) and \(|\Sigma_\beta| \le \mathfrak{c}\) for every \(\beta < \alpha\). Because \(\alpha\) is a countable ordinal, \(\bigl|\bigcup_{\beta<\alpha}\Sigma_\beta\bigr| \le \aleph_0 \cdot \mathfrak{c} = \mathfrak{c}\); hence the collection \(\mathcal{G}\) of sets \(F\) with \(F\) or \(\mathbb{R}^2 \setminus F\) in \(\bigcup_{\beta<\alpha}\Sigma_\beta\) satisfies \(|\mathcal{G}| \le \mathfrak{c}\). Every element of \(\Sigma_\alpha\) is the union of the terms of a sequence in \(\mathcal{G}\), so

\begin{equation*} |\Sigma_\alpha| \le \mathfrak{c}^{\aleph_0} = \bigl(2^{\aleph_0}\bigr)^{\aleph_0} = 2^{\aleph_0} = \mathfrak{c}. \end{equation*}

Thus \(|\Sigma_\alpha| \le \mathfrak{c}\) for every \(\alpha < \omega_1\), and therefore

\begin{equation*} |\mathcal{B}_2| \le |\Sigma| \le \aleph_1 \cdot \mathfrak{c} = \mathfrak{c}, \end{equation*}

where we used \(\aleph_1 \le \mathfrak{c}\).

Step 3: conclusion.

By Cantor’s Theorem \(2^{\mathfrak{c}} > \mathfrak{c} \ge |\mathcal{B}_2|\). Since \(A \mapsto E_A\) is injective on the \(2^{\mathfrak{c}}\) subsets \(A\) of \((0,\infty)\), the sets \(E_A\) cannot all belong to \(\mathcal{B}_2\). Choose \(A \subseteq (0,\infty)\) with \(E_A \notin \mathcal{B}_2\) and set \(E = E_A\). By Step 1, \([E]_a\) and \([E]^a\) are open subsets of \(\mathbb{R}\) for every \(a \in \mathbb{R}\), and \(E \notin \mathcal{B}_2\), as desired.

Two remarks. First, this shows that the converse of 5.6 fails badly: measurability of all cross sections, even openness of all cross sections, does not imply measurability of the set. Second, the argument is a pure counting argument and produces no explicit \(E\); it uses the axiom of choice, both for the cardinal arithmetic in Step 2 and for the comparison \(\aleph_1 \le \mathfrak{c}\).

Problem (5C.3)

Suppose \((X, \mathcal{S})\), \((Y, \mathcal{T})\), and \((Z, \mathcal{U})\) are measurable spaces. We can define \(\mathcal{S} \otimes \mathcal{T} \otimes \mathcal{U}\) to be the smallest \(\sigma\)-algebra on \(X \times Y \times Z\) that contains

\begin{equation*} \{ A \times B \times C : A \in \mathcal{S},\ B \in \mathcal{T},\ C \in \mathcal{U} \}. \end{equation*}

Prove that if we make the obvious identifications of the products \((X \times Y) \times Z\) and \(X \times (Y \times Z)\) with \(X \times Y \times Z\), then

\begin{equation*} \mathcal{S} \otimes \mathcal{T} \otimes \mathcal{U} = (\mathcal{S} \otimes \mathcal{T}) \otimes \mathcal{U} = \mathcal{S} \otimes (\mathcal{T} \otimes \mathcal{U}). \end{equation*}

Solution

Throughout, we identify \(\bigl((x,y),z\bigr)\) and \(\bigl(x,(y,z)\bigr)\) with \((x,y,z)\); under this identification, \((A \times B) \times C\), \(A \times (B \times C)\), and \(A \times B \times C\) are the same subset of \(X \times Y \times Z\), and all three \(\sigma\)-algebras in the statement are collections of subsets of \(X \times Y \times Z\). Write \(\mathcal{P} = \mathcal{S} \otimes \mathcal{T} \otimes \mathcal{U}\), the smallest \(\sigma\)-algebra on \(X \times Y \times Z\) containing every set \(A \times B \times C\) with \(A \in \mathcal{S}\), \(B \in \mathcal{T}\), \(C \in \mathcal{U}\).

First we prove \(\mathcal{P} = (\mathcal{S} \otimes \mathcal{T}) \otimes \mathcal{U}\).

The inclusion \(\mathcal{P} \subseteq (\mathcal{S} \otimes \mathcal{T}) \otimes \mathcal{U}\) is easy. If \(A \in \mathcal{S}\), \(B \in \mathcal{T}\), and \(C \in \mathcal{U}\), then \(A \times B\) is a measurable rectangle in \(\mathcal{S} \otimes \mathcal{T}\), so \(A \times B \in \mathcal{S} \otimes \mathcal{T}\); hence \((A \times B) \times C\) is a measurable rectangle in \((\mathcal{S} \otimes \mathcal{T}) \otimes \mathcal{U}\) and thus lies in \((\mathcal{S} \otimes \mathcal{T}) \otimes \mathcal{U}\). Under our identification \((A \times B) \times C = A \times B \times C\). Thus \((\mathcal{S} \otimes \mathcal{T}) \otimes \mathcal{U}\) is a \(\sigma\)-algebra on \(X \times Y \times Z\) containing all the generators of \(\mathcal{P}\), and therefore contains \(\mathcal{P}\), because \(\mathcal{P}\) is the smallest such \(\sigma\)-algebra.

For the reverse inclusion, it suffices to show that \(D \times C \in \mathcal{P}\) for all \(D \in \mathcal{S} \otimes \mathcal{T}\) and all \(C \in \mathcal{U}\); these sets are the measurable rectangles generating \((\mathcal{S} \otimes \mathcal{T}) \otimes \mathcal{U}\), so once they are known to lie in the \(\sigma\)-algebra \(\mathcal{P}\) we get \((\mathcal{S} \otimes \mathcal{T}) \otimes \mathcal{U} \subseteq \mathcal{P}\).

So temporarily fix \(C \in \mathcal{U}\) and let

\begin{equation*} \mathcal{E} = \{ D \subseteq X \times Y : D \times C \in \mathcal{P} \}. \end{equation*}

We claim \(\mathcal{E}\) is a \(\sigma\)-algebra on \(X \times Y\).

Since \(X \in \mathcal{S}\) and \(Y \in \mathcal{T}\), the set \((X \times Y) \times C = X \times Y \times C\) is one of the generators of \(\mathcal{P}\) and hence lies in \(\mathcal{P}\); thus \(X \times Y \in \mathcal{E}\).

If \(D_1, D_2, \ldots \in \mathcal{E}\), then

\begin{equation*} \Bigl( \bigcup_{k=1}^{\infty} D_k \Bigr) \times C = \bigcup_{k=1}^{\infty} (D_k \times C) \in \mathcal{P}, \end{equation*}

because \(\mathcal{P}\) is closed under countable unions. Hence \(\mathcal{E}\) is closed under countable unions.

If \(D \in \mathcal{E}\), then

\begin{equation*} \bigl( (X \times Y) \setminus D \bigr) \times C = \bigl( (X \times Y \times Z) \setminus (D \times C) \bigr) \cap \bigl( X \times Y \times C \bigr). \end{equation*}

To verify this set equality, note that a point \((x,y,z)\) belongs to the left side if and only if \((x,y) \notin D\) and \(z \in C\); and it belongs to the right side if and only if \(z \in C\) and it is false that both \((x,y) \in D\) and \(z \in C\), which (given \(z \in C\)) says exactly that \((x,y) \notin D\). Since \(D \times C \in \mathcal{P}\) and \(X \times Y \times C \in \mathcal{P}\), and \(\mathcal{P}\) is closed under complementation and finite intersections, the right side lies in \(\mathcal{P}\). Hence \((X \times Y) \setminus D \in \mathcal{E}\), so \(\mathcal{E}\) is closed under complementation.

Thus \(\mathcal{E}\) is a \(\sigma\)-algebra on \(X \times Y\). Moreover if \(A \in \mathcal{S}\) and \(B \in \mathcal{T}\), then \((A \times B) \times C = A \times B \times C \in \mathcal{P}\), so \(A \times B \in \mathcal{E}\). Hence \(\mathcal{E}\) is a \(\sigma\)-algebra containing all the measurable rectangles in \(\mathcal{S} \otimes \mathcal{T}\), which gives \(\mathcal{S} \otimes \mathcal{T} \subseteq \mathcal{E}\). In other words, \(D \times C \in \mathcal{P}\) for every \(D \in \mathcal{S} \otimes \mathcal{T}\). Because \(C \in \mathcal{U}\) was arbitrary, we conclude \((\mathcal{S} \otimes \mathcal{T}) \otimes \mathcal{U} \subseteq \mathcal{P}\), completing the proof that

\begin{equation*} \mathcal{P} = (\mathcal{S} \otimes \mathcal{T}) \otimes \mathcal{U}. \end{equation*}

The proof that \(\mathcal{P} = \mathcal{S} \otimes (\mathcal{T} \otimes \mathcal{U})\) is symmetric. The inclusion \(\mathcal{P} \subseteq \mathcal{S} \otimes (\mathcal{T} \otimes \mathcal{U})\) holds because \(A \times B \times C = A \times (B \times C)\) is a measurable rectangle in \(\mathcal{S} \otimes (\mathcal{T} \otimes \mathcal{U})\) whenever \(A \in \mathcal{S}\), \(B \in \mathcal{T}\), \(C \in \mathcal{U}\). For the reverse inclusion, temporarily fix \(A \in \mathcal{S}\) and let

\begin{equation*} \mathcal{F} = \{ F \subseteq Y \times Z : A \times F \in \mathcal{P} \}. \end{equation*}

Then \(Y \times Z \in \mathcal{F}\) because \(A \times Y \times Z \in \mathcal{P}\); \(\mathcal{F}\) is closed under countable unions because \(A \times \bigcup_k F_k = \bigcup_k (A \times F_k)\); and \(\mathcal{F}\) is closed under complementation because

\begin{equation*} A \times \bigl( (Y \times Z) \setminus F \bigr) = \bigl( (X \times Y \times Z) \setminus (A \times F) \bigr) \cap \bigl( A \times Y \times Z \bigr), \end{equation*}

which is verified exactly as above: \((x,y,z)\) lies in the left side if and only if \(x \in A\) and \((y,z) \notin F\), and this is precisely the condition described by the right side. So \(\mathcal{F}\) is a \(\sigma\)-algebra on \(Y \times Z\), and it contains every measurable rectangle \(B \times C\) with \(B \in \mathcal{T}\), \(C \in \mathcal{U}\) (because \(A \times (B \times C) = A \times B \times C \in \mathcal{P}\)). Hence \(\mathcal{T} \otimes \mathcal{U} \subseteq \mathcal{F}\), which says that \(A \times F \in \mathcal{P}\) for all \(A \in \mathcal{S}\) and \(F \in \mathcal{T} \otimes \mathcal{U}\). These sets generate \(\mathcal{S} \otimes (\mathcal{T} \otimes \mathcal{U})\), so \(\mathcal{S} \otimes (\mathcal{T} \otimes \mathcal{U}) \subseteq \mathcal{P}\).

Combining the two displayed conclusions gives

\begin{equation*} \mathcal{S} \otimes \mathcal{T} \otimes \mathcal{U} = (\mathcal{S} \otimes \mathcal{T}) \otimes \mathcal{U} = \mathcal{S} \otimes (\mathcal{T} \otimes \mathcal{U}), \end{equation*}

as desired. In particular, parentheses may be dispensed with when forming products of three \(\sigma\)-algebras, which is the general version of the associativity noted after 5.39 for the Borel \(\sigma\)-algebras.

Problem (5C.4)

Show that Lebesgue measure on \(\mathbb{R}^n\) is translation invariant. More precisely, show that if \(E \in \mathcal{B}_n\) and \(a \in \mathbb{R}^n\), then \(a + E \in \mathcal{B}_n\) and \(\lambda_n(a + E) = \lambda_n(E)\), where

\begin{equation*} a + E = \{ a + x : x \in E \}. \end{equation*}

Solution

Fix \(a \in \mathbb{R}^n\).

Step 1: \(a + E \in \mathcal{B}_n\) for every \(E \in \mathcal{B}_n\).

Let \(\mathcal{E} = \{ E \in \mathcal{B}_n : a + E \in \mathcal{B}_n \}\). If \(G \subseteq \mathbb{R}^n\) is open, then \(a + G\) is open, because \(a + B(x, \delta) = B(a + x, \delta)\) for all \(x \in \mathbb{R}^n\) and \(\delta > 0\) (translation preserves the \(\|\cdot\|_\infty\) distance), so an open cube centered at \(x\) and contained in \(G\) translates to an open cube centered at \(a + x\) and contained in \(a + G\). Hence \(\mathcal{E}\) contains every open subset of \(\mathbb{R}^n\). Because \(x \mapsto a + x\) is a bijection of \(\mathbb{R}^n\) onto \(\mathbb{R}^n\),

\begin{equation*} a + (\mathbb{R}^n \setminus E) = \mathbb{R}^n \setminus (a + E) \quad \text{and} \quad a + \bigcup_{k=1}^{\infty} E_k = \bigcup_{k=1}^{\infty} (a + E_k), \end{equation*}

so \(\mathcal{E}\) is closed under complementation and countable unions. Thus \(\mathcal{E}\) is a \(\sigma\)-algebra on \(\mathbb{R}^n\) containing all open sets, which gives \(\mathcal{E} = \mathcal{B}_n\). In other words, \(a + E \in \mathcal{B}_n\) whenever \(E \in \mathcal{B}_n\).

Step 2: the case \(n = 1\).

Lebesgue measure \(\lambda_1\) is the restriction of outer measure to the Borel subsets of \(\mathbb{R}\), and outer measure is translation invariant by 2.7. Hence \(\lambda_1(a + E) = |a + E| = |E| = \lambda_1(E)\) for every \(E \in \mathcal{B}_1\) and every \(a \in \mathbb{R}\).

Step 3: translation invariance of the integral in a dimension where translation invariance of the measure is known.

Suppose \(m \in \mathbb{Z}^+\) and we already know that \(\lambda_m(b + F) = \lambda_m(F)\) for all \(F \in \mathcal{B}_m\) and all \(b \in \mathbb{R}^m\). We claim that if \(h : \mathbb{R}^m \to [0, \infty]\) is \(\mathcal{B}_m\)-measurable and \(b \in \mathbb{R}^m\), then the function \(x \mapsto h(x - b)\) is \(\mathcal{B}_m\)-measurable and

\begin{equation*} \int_{\mathbb{R}^m} h(x - b) \, d\lambda_m(x) = \int_{\mathbb{R}^m} h \, d\lambda_m . \end{equation*}

Measurability is immediate from Step 1: if \(B\) is a Borel subset of \([0,\infty]\), then \(\{ x : h(x-b) \in B \} = b + h^{-1}(B) \in \mathcal{B}_m\).

If \(h = \chi_F\) with \(F \in \mathcal{B}_m\), then \(\chi_F(x - b) = 1\) exactly when \(x \in b + F\), so \(x \mapsto h(x-b)\) equals \(\chi_{b+F}\) and

\begin{equation*} \int_{\mathbb{R}^m} \chi_F(x-b) \, d\lambda_m(x) = \lambda_m(b + F) = \lambda_m(F) = \int_{\mathbb{R}^m} \chi_F \, d\lambda_m . \end{equation*}

If \(h = c_1 \chi_{F_1} + \cdots + c_j \chi_{F_j}\) is a nonnegative simple \(\mathcal{B}_m\)-measurable function, then \(h(x - b) = c_1 \chi_{b + F_1}(x) + \cdots + c_j \chi_{b + F_j}(x)\), so the claim follows from the previous sentence together with additivity and homogeneity of the integral of nonnegative functions (3.16 and 3.20).

For general \(\mathcal{B}_m\)-measurable \(h : \mathbb{R}^m \to [0,\infty]\), use 2.89 to choose an increasing sequence \(s_1 \le s_2 \le \cdots\) of nonnegative simple \(\mathcal{B}_m\)-measurable functions with \(\lim_{k \to \infty} s_k(x) = h(x)\) for every \(x\). Then \(k \mapsto s_k(x - b)\) is increasing with limit \(h(x-b)\) for every \(x\). Two applications of the Monotone Convergence Theorem (3.11) give

\begin{equation*} \int_{\mathbb{R}^m} h(x-b) \, d\lambda_m(x) = \lim_{k \to \infty} \int_{\mathbb{R}^m} s_k(x-b) \, d\lambda_m(x) = \lim_{k \to \infty} \int_{\mathbb{R}^m} s_k \, d\lambda_m = \int_{\mathbb{R}^m} h \, d\lambda_m , \end{equation*}

proving the claim.

Step 4: induction on \(n\).

Step 2 gives the result for \(n = 1\). Suppose \(n > 1\) and \(\lambda_{n-1}(b + F) = \lambda_{n-1}(F)\) for all \(F \in \mathcal{B}_{n-1}\) and all \(b \in \mathbb{R}^{n-1}\).

Let \(E \in \mathcal{B}_n\) and write \(a = (a^{\prime}, a_n)\) with \(a^{\prime} \in \mathbb{R}^{n-1}\) and \(a_n \in \mathbb{R}\), thinking of \(\mathbb{R}^n\) as \(\mathbb{R}^{n-1} \times \mathbb{R}\) and of \(\lambda_n\) as \(\lambda_{n-1} \times \lambda_1\) (recall \(\mathcal{B}_n = \mathcal{B}_{n-1} \otimes \mathcal{B}_1\) by 5.39). Every \(\lambda_k\) is a \(\sigma\)-finite measure, because \(\mathbb{R}^k = \bigcup_{m=1}^{\infty} C_m\) where \(C_m\) is the open cube centered at the origin with side length \(m\), and \(\lambda_k(C_m) = m^k < \infty\); so the hypotheses of 5.20 and 5.25 are met throughout what follows. By Step 1, \(a + E \in \mathcal{B}_n\). For \(x \in \mathbb{R}^{n-1}\),

\begin{equation*} [a + E]_x = \{ y \in \mathbb{R} : (x, y) \in a + E \} = \{ y \in \mathbb{R} : (x - a^{\prime}, y - a_n) \in E \} = a_n + [E]_{x - a^{\prime}}, \end{equation*}

where the middle equality holds because \((x,y) \in a + E\) if and only if \((x,y) - a = (x - a^{\prime}, y - a_n) \in E\). Hence, by Step 2,

\begin{equation*} \lambda_1\bigl( [a+E]_x \bigr) = \lambda_1\bigl( a_n + [E]_{x-a^{\prime}} \bigr) = \lambda_1\bigl( [E]_{x - a^{\prime}} \bigr). \end{equation*}

Define \(g : \mathbb{R}^{n-1} \to [0, \infty]\) by \(g(x) = \lambda_1([E]_x)\); this is a \(\mathcal{B}_{n-1}\)-measurable function by 5.20(a), and the equation displayed just above says that \(x \mapsto \lambda_1([a+E]_x)\) is the function \(x \mapsto g(x - a^{\prime})\). Therefore, using the formula for \(\lambda_n\) coming from 5.25,

\begin{equation*} \begin{aligned} \lambda_n(a + E) &= \int_{\mathbb{R}^{n-1}} \lambda_1\bigl( [a+E]_x \bigr) \, d\lambda_{n-1}(x) \\ &= \int_{\mathbb{R}^{n-1}} g(x - a^{\prime}) \, d\lambda_{n-1}(x) \\ &= \int_{\mathbb{R}^{n-1}} g \, d\lambda_{n-1} \\ &= \int_{\mathbb{R}^{n-1}} \lambda_1\bigl( [E]_x \bigr) \, d\lambda_{n-1}(x) \\ &= \lambda_n(E), \end{aligned} \end{equation*}

where the third equality is Step 3 applied with \(m = n - 1\) and \(b = a^{\prime}\) (legitimate by the induction hypothesis).

This completes the induction, so \(\lambda_n(a + E) = \lambda_n(E)\) for every \(n \in \mathbb{Z}^+\), every \(E \in \mathcal{B}_n\), and every \(a \in \mathbb{R}^n\). Together with Step 1, Lebesgue measure on \(\mathbb{R}^n\) is translation invariant.

Problem (5C.5)

Suppose \(f : \mathbb{R}^n \to \mathbb{R}\) is \(\mathcal{B}_n\)-measurable and \(t \in \mathbb{R} \setminus \{0\}\). Define \(f_t : \mathbb{R}^n \to \mathbb{R}\) by \(f_t(x) = f(tx)\).

(a) Prove that \(f_t\) is \(\mathcal{B}_n\)-measurable.

(b) Prove that if \(\int_{\mathbb{R}^n} f \, d\lambda_n\) is defined, then

\begin{equation*} \int_{\mathbb{R}^n} f_t \, d\lambda_n = \frac{1}{|t|^n} \int_{\mathbb{R}^n} f \, d\lambda_n . \end{equation*}

Solution

We first extend 5.41, which is stated for \(t > 0\), to all nonzero \(t\).

Lemma. If \(t \in \mathbb{R} \setminus \{0\}\) and \(E \in \mathcal{B}_n\), then \(tE \in \mathcal{B}_n\) and \(\lambda_n(tE) = |t|^n \lambda_n(E)\).

Proof of the Lemma. For \(t > 0\) this is exactly 5.41. So it suffices to treat the map \(x \mapsto -x\), since for \(t < 0\) we have \(tE = -\bigl(|t| E\bigr)\).

Write \(-F = \{ -x : x \in F \}\). The collection \(\{ F \in \mathcal{B}_n : -F \in \mathcal{B}_n \}\) contains every open set (if \(G\) is open then so is \(-G\), because \(-B(x,\delta) = B(-x,\delta)\)) and is closed under complementation and countable unions (because \(x \mapsto -x\) is a bijection of \(\mathbb{R}^n\), so \(-(\mathbb{R}^n \setminus F) = \mathbb{R}^n \setminus (-F)\) and \(-\bigcup_k F_k = \bigcup_k (-F_k)\)). Hence it is a \(\sigma\)-algebra containing the open sets, so it equals \(\mathcal{B}_n\): we get \(-F \in \mathcal{B}_n\) for all \(F \in \mathcal{B}_n\).

Next, \(\lambda_n(-F) = \lambda_n(F)\) for all \(F \in \mathcal{B}_n\), by induction on \(n\). For \(n = 1\): \(\lambda_1\) is the restriction of outer measure, and if \(I_1, I_2, \ldots\) is a sequence of open intervals whose union contains \(F\), then \(-I_1, -I_2, \ldots\) is a sequence of open intervals of the same lengths whose union contains \(-F\) (and conversely); hence \(|-F| = |F|\), so \(\lambda_1(-F) = \lambda_1(F)\).

Now let \(n > 1\) and assume \(\lambda_{n-1}(-F) = \lambda_{n-1}(F)\) for all \(F \in \mathcal{B}_{n-1}\). Exactly as in Step 3 of the solution to Exercise 4 in this section (first characteristic functions, then nonnegative simple functions using 3.16 and 3.20, then general nonnegative measurable functions using 2.89 and the Monotone Convergence Theorem 3.11), the induction hypothesis implies that

\begin{equation*} \int_{\mathbb{R}^{n-1}} h(-x) \, d\lambda_{n-1}(x) = \int_{\mathbb{R}^{n-1}} h \, d\lambda_{n-1} \end{equation*}

for every \(\mathcal{B}_{n-1}\)-measurable \(h : \mathbb{R}^{n-1} \to [0, \infty]\); note that \(x \mapsto h(-x)\) is \(\mathcal{B}_{n-1}\)-measurable because \(\{ x : h(-x) \in B \} = -h^{-1}(B) \in \mathcal{B}_{n-1}\), and that \(\chi_F(-x) = \chi_{-F}(x)\).

Let \(F \in \mathcal{B}_n\) and view \(\mathbb{R}^n = \mathbb{R}^{n-1} \times \mathbb{R}\), so \(\lambda_n = \lambda_{n-1} \times \lambda_1\) and \(\mathcal{B}_n = \mathcal{B}_{n-1} \otimes \mathcal{B}_1\) (5.39). Each \(\lambda_k\) is \(\sigma\)-finite, since \(\mathbb{R}^k\) is the union of the increasing sequence of open cubes centered at the origin with side lengths \(1, 2, 3, \ldots\), each of finite measure; so 5.20 and 5.25 apply. For \(x \in \mathbb{R}^{n-1}\),

\begin{equation*} [-F]_x = \{ y \in \mathbb{R} : (x,y) \in -F \} = \{ y \in \mathbb{R} : (-x, -y) \in F \} = -[F]_{-x} . \end{equation*}

Hence \(\lambda_1([-F]_x) = \lambda_1([F]_{-x}) = g(-x)\), where \(g(x) = \lambda_1([F]_x)\) is \(\mathcal{B}_{n-1}\)-measurable by 5.20(a). Therefore, using 5.25,

\begin{equation*} \lambda_n(-F) = \int_{\mathbb{R}^{n-1}} g(-x) \, d\lambda_{n-1}(x) = \int_{\mathbb{R}^{n-1}} g \, d\lambda_{n-1} = \lambda_n(F). \end{equation*}

This completes the induction.

Finally, for \(t < 0\) and \(E \in \mathcal{B}_n\) we have \(|t| E \in \mathcal{B}_n\) by 5.41, hence \(tE = -(|t|E) \in \mathcal{B}_n\), and

\begin{equation*} \lambda_n(tE) = \lambda_n\bigl( -(|t|E) \bigr) = \lambda_n\bigl( |t| E \bigr) = |t|^n \lambda_n(E), \end{equation*}

where the last equality is 5.41. This proves the Lemma.

(a) Let \(B \subseteq \mathbb{R}\) be a Borel set. Since \(t \ne 0\),

\begin{equation*} f_t^{-1}(B) = \{ x \in \mathbb{R}^n : f(tx) \in B \} = \{ x \in \mathbb{R}^n : tx \in f^{-1}(B) \} = \tfrac{1}{t} f^{-1}(B). \end{equation*}

Because \(f\) is \(\mathcal{B}_n\)-measurable, \(f^{-1}(B) \in \mathcal{B}_n\), and then the Lemma (applied with the nonzero number \(1/t\)) gives \(\tfrac{1}{t} f^{-1}(B) \in \mathcal{B}_n\). Thus \(f_t^{-1}(B) \in \mathcal{B}_n\) for every Borel \(B \subseteq \mathbb{R}\), which by 2.35 says that \(f_t\) is \(\mathcal{B}_n\)-measurable.

(b) We build up in the usual way. Throughout, for a function \(h\) on \(\mathbb{R}^n\) we write \(h_t(x) = h(tx)\); part (a) and its proof show that \(h_t\) is \(\mathcal{B}_n\)-measurable whenever \(h\) is (the argument in (a) never used that \(f\) is real valued, so it applies to \(\mathcal{B}_n\)-measurable functions with values in \([0,\infty]\) as well).

Characteristic functions. Let \(E \in \mathcal{B}_n\) and \(h = \chi_E\). Then \(h_t(x) = \chi_E(tx)\), which equals \(1\) exactly when \(tx \in E\), that is, exactly when \(x \in \tfrac{1}{t}E\). Thus \(h_t = \chi_{(1/t)E}\) and, by the Lemma,

\begin{equation*} \int_{\mathbb{R}^n} h_t \, d\lambda_n = \lambda_n\Bigl( \tfrac{1}{t} E \Bigr) = \Bigl| \tfrac{1}{t} \Bigr|^{\,n} \lambda_n(E) = \frac{1}{|t|^n} \int_{\mathbb{R}^n} h \, d\lambda_n . \end{equation*}

Nonnegative simple functions. If \(h = c_1 \chi_{E_1} + \cdots + c_j \chi_{E_j}\) with \(c_1, \ldots, c_j \in [0,\infty)\) and \(E_1, \ldots, E_j \in \mathcal{B}_n\), then \(h_t = c_1 (\chi_{E_1})_t + \cdots + c_j (\chi_{E_j})_t\), so additivity and homogeneity of the integral of nonnegative functions (3.16 and 3.20) together with the previous paragraph give \(\int h_t \, d\lambda_n = |t|^{-n} \int h \, d\lambda_n\).

Nonnegative measurable functions. Let \(h : \mathbb{R}^n \to [0,\infty]\) be \(\mathcal{B}_n\)-measurable. By 2.89 there is an increasing sequence \(s_1 \le s_2 \le \cdots\) of nonnegative simple \(\mathcal{B}_n\)-measurable functions converging pointwise to \(h\). Then \((s_1)_t \le (s_2)_t \le \cdots\) converges pointwise to \(h_t\) (evaluate at \(tx\)). Applying the Monotone Convergence Theorem (3.11) twice,

\begin{equation*} \int_{\mathbb{R}^n} h_t \, d\lambda_n = \lim_{k \to \infty} \int_{\mathbb{R}^n} (s_k)_t \, d\lambda_n = \lim_{k \to \infty} \frac{1}{|t|^n} \int_{\mathbb{R}^n} s_k \, d\lambda_n = \frac{1}{|t|^n} \int_{\mathbb{R}^n} h \, d\lambda_n . \end{equation*}

The general case. Suppose \(\int_{\mathbb{R}^n} f \, d\lambda_n\) is defined; by 3.18 this means that \(f\) is \(\mathcal{B}_n\)-measurable and at least one of \(\int f^+ \, d\lambda_n\) and \(\int f^- \, d\lambda_n\) is finite. For every \(x \in \mathbb{R}^n\),

\begin{equation*} (f_t)^+(x) = \max\{ f(tx), 0 \} = f^+(tx) = (f^+)_t(x), \qquad (f_t)^-(x) = (f^-)_t(x), \end{equation*}

so \((f_t)^{\pm} = (f^{\pm})_t\). By the previous paragraph,

\begin{equation*} \int_{\mathbb{R}^n} (f_t)^+ \, d\lambda_n = \frac{1}{|t|^n} \int_{\mathbb{R}^n} f^+ \, d\lambda_n, \qquad \int_{\mathbb{R}^n} (f_t)^- \, d\lambda_n = \frac{1}{|t|^n} \int_{\mathbb{R}^n} f^- \, d\lambda_n . \end{equation*}

Because \(0 < |t|^{-n} < \infty\), at least one of the two integrals on the left is finite; hence \(\int_{\mathbb{R}^n} f_t \, d\lambda_n\) is defined and

\begin{equation*} \int_{\mathbb{R}^n} f_t \, d\lambda_n = \int_{\mathbb{R}^n} (f_t)^+ \, d\lambda_n - \int_{\mathbb{R}^n} (f_t)^- \, d\lambda_n = \frac{1}{|t|^n} \Bigl( \int_{\mathbb{R}^n} f^+ \, d\lambda_n - \int_{\mathbb{R}^n} f^- \, d\lambda_n \Bigr) = \frac{1}{|t|^n} \int_{\mathbb{R}^n} f \, d\lambda_n , \end{equation*}

where the subtraction is legitimate because not both terms are \(\infty\). This completes the proof of (b).

Problem (5C.6)

Suppose \(\lambda\) denotes Lebesgue measure on \((\mathbb{R}, \mathcal{L})\), where \(\mathcal{L}\) is the \(\sigma\)-algebra of Lebesgue measurable subsets of \(\mathbb{R}\). Show that there exist subsets \(E\) and \(F\) of \(\mathbb{R}^2\) such that

  • \(F \in \mathcal{L} \otimes \mathcal{L}\) and \((\lambda \times \lambda)(F) = 0\);
  • \(E \subseteq F\) but \(E \notin \mathcal{L} \otimes \mathcal{L}\).

[The measure space \((\mathbb{R}, \mathcal{L}, \lambda)\) has the property that every subset of a measurable set with measure \(0\) is measurable. This exercise asks you to show that the measure space \((\mathbb{R}^2, \mathcal{L} \otimes \mathcal{L}, \lambda \times \lambda)\) does not have this property.]

Solution

First note that \(\lambda\) on \((\mathbb{R}, \mathcal{L})\) is a \(\sigma\)-finite measure, because \(\mathbb{R} = \bigcup_{k=1}^{\infty} (-k, k)\) and \(\lambda\bigl( (-k,k) \bigr) = 2k < \infty\) for each \(k \in \mathbb{Z}^+\). Hence the product measure \(\lambda \times \lambda\) on \((\mathbb{R}^2, \mathcal{L} \otimes \mathcal{L})\) is defined (5.25) and is a measure (5.27).

There exists a set \(V \subseteq \mathbb{R}\) that is not Lebesgue measurable. Indeed, 2.18 produces disjoint sets \(A, B \subseteq \mathbb{R}\) with \(|A \cup B| \ne |A| + |B|\), whereas outer measure is a measure on \((\mathbb{R}, \mathcal{L})\) by 2.72(b) and hence is finitely additive on disjoint sets in \(\mathcal{L}\); so \(A\) and \(B\) cannot both lie in \(\mathcal{L}\). (This is the observation made in the discussion following 2.72.) Fix such a \(V\), so \(V \notin \mathcal{L}\).

Define

\begin{equation*} F = \{0\} \times \mathbb{R} \qquad \text{and} \qquad E = \{0\} \times V . \end{equation*}

Clearly \(E \subseteq F\).

\(F\) is a measurable rectangle: \(\{0\} \in \mathcal{L}\) (indeed \(\{0\}\) is a Borel set, and every Borel set is Lebesgue measurable) and \(\mathbb{R} \in \mathcal{L}\). Hence \(F \in \mathcal{L} \otimes \mathcal{L}\).

Next we compute \((\lambda \times \lambda)(F)\). For \(x \in \mathbb{R}\) the cross section is

\begin{equation*} [F]_x = \begin{cases} \mathbb{R} & \text{if } x = 0, \\ \emptyset & \text{if } x \ne 0, \end{cases} \end{equation*}

so \(\lambda([F]_x) = \infty \cdot \chi_{\{0\}}(x)\). By 5.25,

\begin{equation*} (\lambda \times \lambda)(F) = \int_{\mathbb{R}} \int_{\mathbb{R}} \chi_F(x,y) \, d\lambda(y) \, d\lambda(x) = \int_{\mathbb{R}} \lambda\bigl( [F]_x \bigr) \, d\lambda(x) = \int_{\mathbb{R}} \infty \cdot \chi_{\{0\}} \, d\lambda = 0, \end{equation*}

the last equality holding because the integrand vanishes outside the set \(\{0\}\), which has \(\lambda\)-measure \(0\): for every \(c \in [0,\infty)\) we have \(\int c \chi_{\{0\}} \, d\lambda = c \lambda(\{0\}) = 0\), and \(\infty \cdot \chi_{\{0\}}\) is the increasing pointwise limit of \(k \chi_{\{0\}}\) as \(k \to \infty\), so the Monotone Convergence Theorem (3.11) gives \(\int \infty \cdot \chi_{\{0\}} \, d\lambda = \lim_{k \to \infty} 0 = 0\). (Equivalently, use 5.26: \((\lambda \times \lambda)(F) = \lambda(\{0\}) \cdot \lambda(\mathbb{R}) = 0 \cdot \infty = 0\).) Thus \(F\) satisfies the first bullet point.

Finally, suppose toward a contradiction that \(E \in \mathcal{L} \otimes \mathcal{L}\). By 5.6, every cross section of a set in \(\mathcal{L} \otimes \mathcal{L}\) is measurable; in particular \([E]_0 \in \mathcal{L}\). But

\begin{equation*} [E]_0 = \{ y \in \mathbb{R} : (0, y) \in E \} = V, \end{equation*}

so \(V \in \mathcal{L}\), contradicting the choice of \(V\). Hence \(E \notin \mathcal{L} \otimes \mathcal{L}\), and \(E\) satisfies the second bullet point.

Thus \(E \subseteq F\), \(F \in \mathcal{L} \otimes \mathcal{L}\) with \((\lambda \times \lambda)(F) = 0\), and yet \(E \notin \mathcal{L} \otimes \mathcal{L}\). In the terminology of the bracketed remark, the measure space \((\mathbb{R}^2, \mathcal{L} \otimes \mathcal{L}, \lambda \times \lambda)\) is not complete, even though \((\mathbb{R}, \mathcal{L}, \lambda)\) is.

Problem (5C.7)

Suppose \(m \in \mathbb{Z}^+\). Verify that the collection of sets \(\mathcal{E}_m\) that appears in the proof of 5.41 is a monotone class.

Solution

Recall the setting of the proof of 5.41: a number \(t > 0\) is fixed, \(C_m\) denotes the open cube in \(\mathbb{R}^n\) centered at the origin with side length \(m\), and

\begin{equation*} \mathcal{E}_m = \{ E \in \mathcal{B}_n : E \subseteq C_m \text{ and } \lambda_n(tE) = t^n \lambda_n(E) \}. \end{equation*}

Recall also from the first paragraph of that proof that \(tE \in \mathcal{B}_n\) for every \(E \in \mathcal{B}_n\), so the condition \(\lambda_n(tE) = t^n \lambda_n(E)\) makes sense. By 5.15, we must show that \(\mathcal{E}_m\) is closed under countable increasing unions and under countable decreasing intersections.

Two elementary observations will be used repeatedly. First, because \(t > 0\), the map \(x \mapsto tx\) is a bijection of \(\mathbb{R}^n\) onto \(\mathbb{R}^n\), and therefore it carries unions to unions and intersections to intersections:

\begin{equation*} t \Bigl( \bigcup_{k=1}^{\infty} E_k \Bigr) = \bigcup_{k=1}^{\infty} (t E_k) \qquad \text{and} \qquad t \Bigl( \bigcap_{k=1}^{\infty} E_k \Bigr) = \bigcap_{k=1}^{\infty} (t E_k). \end{equation*}

(The inclusion \(\subseteq\) in the second equation holds for any map; the reverse inclusion uses injectivity of \(x \mapsto tx\).) Also, \(E \subseteq F\) implies \(tE \subseteq tF\).

Second, \(\lambda_n(C_m) = m^n < \infty\). Indeed \(C_m = (-\tfrac{m}{2}, \tfrac{m}{2})^n\), and by induction on \(n\), using \(\lambda_n = \lambda_{n-1} \times \lambda_1\) and 5.26,

\begin{equation*} \lambda_n(C_m) = \lambda_{n-1}\bigl( (-\tfrac{m}{2}, \tfrac{m}{2})^{n-1} \bigr) \cdot \lambda_1\bigl( (-\tfrac{m}{2}, \tfrac{m}{2}) \bigr) = m^{n-1} \cdot m = m^n , \end{equation*}

with the case \(n = 1\) holding because \(\lambda_1\) of an open interval is its length. Similarly \(tC_m\) is the open cube centered at the origin with side length \(tm\), so \(\lambda_n(t C_m) = ( tm)^n < \infty\).

Countable increasing unions. Suppose \(E_1 \subseteq E_2 \subseteq \cdots\) is an increasing sequence of sets in \(\mathcal{E}_m\), and let \(E = \bigcup_{k=1}^{\infty} E_k\). Then \(E \in \mathcal{B}_n\) (as \(\mathcal{B}_n\) is a \(\sigma\)-algebra) and \(E \subseteq C_m\) (as each \(E_k \subseteq C_m\)). Moreover \(tE_1 \subseteq tE_2 \subseteq \cdots\) is an increasing sequence in \(\mathcal{B}_n\) with union \(tE\). Hence, using 2.59 twice,

\begin{equation*} \lambda_n(tE) = \lim_{k \to \infty} \lambda_n(t E_k) = \lim_{k \to \infty} t^n \lambda_n(E_k) = t^n \lim_{k \to \infty} \lambda_n(E_k) = t^n \lambda_n(E), \end{equation*}

where the second equality holds because each \(E_k \in \mathcal{E}_m\). Therefore \(E \in \mathcal{E}_m\).

Countable decreasing intersections. Suppose \(E_1 \supseteq E_2 \supseteq \cdots\) is a decreasing sequence of sets in \(\mathcal{E}_m\), and let \(E = \bigcap_{k=1}^{\infty} E_k\). Then \(E \in \mathcal{B}_n\) and \(E \subseteq C_m\). Also \(tE_1 \supseteq tE_2 \supseteq \cdots\) is a decreasing sequence in \(\mathcal{B}_n\) with intersection \(tE\).

The finiteness hypothesis needed for 2.60 holds in both of the sequences we use: since \(E_1 \subseteq C_m\), we have \(\lambda_n(E_1) \le \lambda_n(C_m) = m^n < \infty\) by 2.57(a); and \(tE_1 \subseteq tC_m\), so \(\lambda_n(tE_1) \le \lambda_n(tC_m) = ( tm)^n < \infty\). This is exactly why the sets in \(\mathcal{E}_m\) are required to sit inside the bounded cube \(C_m\).

Hence, applying 2.60 to \(tE_1 \supseteq tE_2 \supseteq \cdots\) and then to \(E_1 \supseteq E_2 \supseteq \cdots\),

\begin{equation*} \lambda_n(tE) = \lim_{k \to \infty} \lambda_n(t E_k) = \lim_{k \to \infty} t^n \lambda_n(E_k) = t^n \lim_{k \to \infty} \lambda_n(E_k) = t^n \lambda_n(E), \end{equation*}

where again the second equality uses \(E_k \in \mathcal{E}_m\). Therefore \(E \in \mathcal{E}_m\).

We have shown that \(\mathcal{E}_m\) is closed under countable increasing unions and under countable decreasing intersections, so \(\mathcal{E}_m\) is a monotone class on \(C_m\), as claimed in the proof of 5.41.

Problem (5C.8)

Show that the open unit ball in \(\mathbb{R}^n\) is an open subset of \(\mathbb{R}^n\).

Solution

Recall (5.43) that the open unit ball is

\begin{equation*} B_n = \{(x_1,\dots,x_n) \in \mathbb{R}^n : x_1{}^2 + \cdots + x_n{}^2 < 1\}, \end{equation*}

and that a subset \(G\) of \(\mathbb{R}^n\) is open if for every \(x \in G\) there exists \(\delta > 0\) with \(B(x,\delta) \subseteq G\), where \(B(x,\delta) = \{y \in \mathbb{R}^n : \|y - x\|_\infty < \delta\}\) and \(\|(y_1,\dots,y_n)\|_\infty = \max\{|y_1|,\dots,|y_n|\}\).

Write \(\|y\|_2 = (y_1{}^2 + \cdots + y_n{}^2)^{1/2}\), so that \(B_n = \{y \in \mathbb{R}^n : \|y\|_2 < 1\}\). Two standard facts about \(\|\cdot\|_2\) will be used.

First, the triangle inequality \(\|u + v\|_2 \le \|u\|_2 + \|v\|_2\). This follows from the Cauchy–Schwarz inequality \(\sum_{j=1}^n u_j v_j \le \|u\|_2 \|v\|_2\), because

\begin{equation*} \|u+v\|_2{}^2 = \|u\|_2{}^2 + 2\sum_{j=1}^n u_j v_j + \|v\|_2{}^2 \le \|u\|_2{}^2 + 2\|u\|_2\|v\|_2 + \|v\|_2{}^2 = (\|u\|_2 + \|v\|_2)^2 . \end{equation*}

Second, the comparison

\begin{equation*} \|y\|_2 = \Bigl(\sum_{j=1}^n y_j{}^2\Bigr)^{1/2} \le \Bigl(\sum_{j=1}^n \|y\|_\infty{}^2\Bigr)^{1/2} = \sqrt{n}\,\|y\|_\infty , \end{equation*}

which holds because \(|y_j| \le \|y\|_\infty\) for each \(j\).

Now suppose \(x \in B_n\). Then \(\|x\|_2 < 1\), so

\begin{equation*} r = 1 - \|x\|_2 > 0 . \end{equation*}

Set \(\delta = r/\sqrt{n} > 0\). Suppose \(y \in B(x,\delta)\), meaning \(\|y - x\|_\infty < \delta\). Then

\begin{equation*} \|y - x\|_2 \le \sqrt{n}\,\|y - x\|_\infty < \sqrt{n}\,\delta = r , \end{equation*}

and hence, by the triangle inequality applied to \(y = x + (y - x)\),

\begin{equation*} \|y\|_2 \le \|x\|_2 + \|y - x\|_2 < \|x\|_2 + r = 1 . \end{equation*}

Thus \(y \in B_n\). Since \(y\) was an arbitrary element of \(B(x,\delta)\), we have shown \(B(x,\delta) \subseteq B_n\).

Every point of \(B_n\) therefore lies in an open cube contained in \(B_n\), which means \(B_n\) is an open subset of \(\mathbb{R}^n\). In particular \(B_n\) is a Borel subset of \(\mathbb{R}^n\), as asserted after 5.43.

Problem (5C.9)

Suppose \(G_1\) is a nonempty subset of \(\mathbb{R}^m\) and \(G_2\) is a nonempty subset of \(\mathbb{R}^n\). Prove that \(G_1 \times G_2\) is an open subset of \(\mathbb{R}^m \times \mathbb{R}^n\) if and only if \(G_1\) is an open subset of \(\mathbb{R}^m\) and \(G_2\) is an open subset of \(\mathbb{R}^n\).

[One direction of this result was already proved (see 5.36); both directions are stated here to make the result look prettier and to be comparable to the next exercise, where neither direction has been proved.]

Solution

Throughout we use the identification of \(\mathbb{R}^m \times \mathbb{R}^n\) with \(\mathbb{R}^{m+n}\) and the identity

\begin{equation*} B(x,\delta) \times B(y,\delta) = B\bigl((x,y),\delta\bigr) \end{equation*}

for \(x \in \mathbb{R}^m\), \(y \in \mathbb{R}^n\), \(\delta > 0\), both recorded in the review at the start of Section 5C. This identity holds because \(\|(u,v)\|_\infty = \max\{\|u\|_\infty, \|v\|_\infty\}\) under that identification.

First suppose \(G_1\) is open in \(\mathbb{R}^m\) and \(G_2\) is open in \(\mathbb{R}^n\). Then \(G_1 \times G_2\) is open in \(\mathbb{R}^{m+n}\) by 5.36. (This direction does not use the hypothesis that \(G_1\) and \(G_2\) are nonempty.)

Now suppose \(G_1 \times G_2\) is an open subset of \(\mathbb{R}^{m+n}\). Because \(G_2 \neq \varnothing\), we may fix some \(y_0 \in G_2\).

To prove \(G_1\) is open, let \(x \in G_1\). Then \((x,y_0) \in G_1 \times G_2\), so openness of \(G_1 \times G_2\) gives \(\delta > 0\) with

\begin{equation*} B\bigl((x,y_0),\delta\bigr) \subseteq G_1 \times G_2 . \end{equation*}

Let \(u \in B(x,\delta)\). Since \(y_0 \in B(y_0,\delta)\), we get

\begin{equation*} (u,y_0) \in B(x,\delta) \times B(y_0,\delta) = B\bigl((x,y_0),\delta\bigr) \subseteq G_1 \times G_2 , \end{equation*}

and therefore \(u \in G_1\). Thus \(B(x,\delta) \subseteq G_1\). As \(x \in G_1\) was arbitrary, \(G_1\) is an open subset of \(\mathbb{R}^m\).

The proof that \(G_2\) is open is symmetric: because \(G_1 \neq \varnothing\) we may fix \(x_0 \in G_1\); given \(y \in G_2\), choose \(\delta > 0\) with \(B((x_0,y),\delta) \subseteq G_1 \times G_2\), and then for each \(v \in B(y,\delta)\) we have \((x_0,v) \in B(x_0,\delta) \times B(y,\delta) = B((x_0,y),\delta) \subseteq G_1 \times G_2\), so \(v \in G_2\). Hence \(B(y,\delta) \subseteq G_2\) and \(G_2\) is open in \(\mathbb{R}^n\).

The hypothesis that both sets are nonempty is genuinely needed for the second direction: if \(G_2 = \varnothing\) then \(G_1 \times G_2 = \varnothing\) is open no matter what \(G_1\) is, so no conclusion about \(G_1\) could be drawn.

Problem (5C.10)

Suppose \(F_1\) is a nonempty subset of \(\mathbb{R}^m\) and \(F_2\) is a nonempty subset of \(\mathbb{R}^n\). Prove that \(F_1 \times F_2\) is a closed subset of \(\mathbb{R}^m \times \mathbb{R}^n\) if and only if \(F_1\) is a closed subset of \(\mathbb{R}^m\) and \(F_2\) is a closed subset of \(\mathbb{R}^n\).

Solution

As in the previous exercise we identify \(\mathbb{R}^m \times \mathbb{R}^n\) with \(\mathbb{R}^{m+n}\) and use \(B(x,\delta) \times B(y,\delta) = B((x,y),\delta)\). Recall that a subset of \(\mathbb{R}^k\) is closed exactly when its complement is open.

First suppose \(F_1\) is closed in \(\mathbb{R}^m\) and \(F_2\) is closed in \(\mathbb{R}^n\). A point \((x,y)\) fails to lie in \(F_1 \times F_2\) precisely when \(x \notin F_1\) or \(y \notin F_2\). Hence

\begin{equation*} \mathbb{R}^{m+n} \setminus (F_1 \times F_2) = \bigl[(\mathbb{R}^m \setminus F_1) \times \mathbb{R}^n\bigr] \cup \bigl[\mathbb{R}^m \times (\mathbb{R}^n \setminus F_2)\bigr]. \end{equation*}

The sets \(\mathbb{R}^m \setminus F_1\) and \(\mathbb{R}^n \setminus F_2\) are open by hypothesis, and \(\mathbb{R}^m\) and \(\mathbb{R}^n\) are open in themselves. Thus each of the two sets on the right side above is a product of an open subset of \(\mathbb{R}^m\) with an open subset of \(\mathbb{R}^n\), hence is open in \(\mathbb{R}^{m+n}\) by 5.36. A union of two open sets is open. Therefore \(\mathbb{R}^{m+n} \setminus (F_1 \times F_2)\) is open, which means \(F_1 \times F_2\) is closed. (This direction does not use nonemptiness.)

Now suppose \(F_1 \times F_2\) is a closed subset of \(\mathbb{R}^{m+n}\); thus \(\mathbb{R}^{m+n} \setminus (F_1 \times F_2)\) is open. Because \(F_2 \neq \varnothing\), fix \(y_0 \in F_2\).

To prove \(F_1\) is closed we show \(\mathbb{R}^m \setminus F_1\) is open. Let \(x \in \mathbb{R}^m \setminus F_1\). Then \((x,y_0) \notin F_1 \times F_2\), so there exists \(\delta > 0\) with

\begin{equation*} B\bigl((x,y_0),\delta\bigr) \subseteq \mathbb{R}^{m+n} \setminus (F_1 \times F_2). \end{equation*}

Let \(u \in B(x,\delta)\). Then \((u,y_0) \in B(x,\delta) \times B(y_0,\delta) = B((x,y_0),\delta)\), so \((u,y_0) \notin F_1 \times F_2\). Since \(y_0 \in F_2\), this forces \(u \notin F_1\). Hence \(B(x,\delta) \subseteq \mathbb{R}^m \setminus F_1\). As \(x\) was an arbitrary point of \(\mathbb{R}^m \setminus F_1\), that complement is open and so \(F_1\) is closed in \(\mathbb{R}^m\).

The argument that \(F_2\) is closed is symmetric, using a fixed \(x_0 \in F_1\) (available since \(F_1 \neq \varnothing\)): given \(y \in \mathbb{R}^n \setminus F_2\), choose \(\delta > 0\) with \(B((x_0,y),\delta) \cap (F_1 \times F_2) = \varnothing\); for \(v \in B(y,\delta)\) we have \((x_0,v) \in B((x_0,y),\delta)\), so \((x_0,v) \notin F_1 \times F_2\), and since \(x_0 \in F_1\) this gives \(v \notin F_2\). Hence \(B(y,\delta) \subseteq \mathbb{R}^n \setminus F_2\), proving \(F_2\) closed.

Again nonemptiness is needed only for the second direction: \(\varnothing \times F_2 = \varnothing\) is closed regardless of \(F_2\).

Problem (5C.11)

Suppose \(E\) is a subset of \(\mathbb{R}^m \times \mathbb{R}^n\) and

\begin{equation*} A = \{x \in \mathbb{R}^m : (x,y) \in E \text{ for some } y \in \mathbb{R}^n\}. \end{equation*}

(a) Prove that if \(E\) is an open subset of \(\mathbb{R}^m \times \mathbb{R}^n\), then \(A\) is an open subset of \(\mathbb{R}^m\).

(b) Prove or give a counterexample: If \(E\) is a closed subset of \(\mathbb{R}^m \times \mathbb{R}^n\), then \(A\) is a closed subset of \(\mathbb{R}^m\).

Solution

The set \(A\) is the projection of \(E\) onto the first factor. As before we identify \(\mathbb{R}^m \times \mathbb{R}^n\) with \(\mathbb{R}^{m+n}\) and use \(B(x,\delta) \times B(y,\delta) = B((x,y),\delta)\).

(a) Suppose \(E\) is open and let \(x \in A\). By the definition of \(A\) there exists \(y \in \mathbb{R}^n\) with \((x,y) \in E\). Because \(E\) is open, there exists \(\delta > 0\) with

\begin{equation*} B\bigl((x,y),\delta\bigr) = B(x,\delta) \times B(y,\delta) \subseteq E . \end{equation*}

Now let \(u \in B(x,\delta)\). Since \(y \in B(y,\delta)\), we have \((u,y) \in B(x,\delta) \times B(y,\delta) \subseteq E\). So \(y\) witnesses that \(u \in A\). Hence \(B(x,\delta) \subseteq A\).

Every point of \(A\) is thus the center of an open cube contained in \(A\), so \(A\) is an open subset of \(\mathbb{R}^m\).

(b) The statement is false; here is a counterexample. Take \(m = n = 1\) and

\begin{equation*} E = \{(x,y) \in \mathbb{R}^2 : xy = 1\}, \end{equation*}

the hyperbola.

\(E\) is closed. To see this directly, suppose \((x_0,y_0) \notin E\), so \(x_0 y_0 \neq 1\). Put \(\varepsilon = |x_0 y_0 - 1| > 0\) and choose \(\delta > 0\) small enough that \(\delta \le 1\) and \(\delta(|x_0| + |y_0| + 1) < \varepsilon\). If \(\|(x,y) - (x_0,y_0)\|_\infty < \delta\), then \(|x - x_0| < \delta\), \(|y - y_0| < \delta\), and \(|y| < |y_0| + \delta \le |y_0| + 1\), so

\begin{equation*} |xy - x_0 y_0| \le |y|\,|x - x_0| + |x_0|\,|y - y_0| < (|y_0| + 1)\delta + |x_0|\delta = \delta(|x_0| + |y_0| + 1) < \varepsilon . \end{equation*}

Hence \(xy \neq 1\), i.e. \((x,y) \notin E\). Thus \(B((x_0,y_0),\delta) \subseteq \mathbb{R}^2 \setminus E\), so \(\mathbb{R}^2 \setminus E\) is open and \(E\) is closed. (Equivalently: \(E\) is the inverse image of the closed set \(\{1\}\) under the continuous map \((x,y) \mapsto xy\).)

Now compute \(A\). If \(x \neq 0\) then \((x, 1/x) \in E\), so \(x \in A\). If \(x = 0\) then \(xy = 0 \neq 1\) for every \(y\), so \(0 \notin A\). Therefore

\begin{equation*} A = \mathbb{R} \setminus \{0\}, \end{equation*}

whose complement \(\{0\}\) is not open in \(\mathbb{R}\) (every open cube centered at \(0\) contains nonzero points). Hence \(A\) is not closed.

The same phenomenon occurs for arbitrary \(m\) and \(n\): take

\begin{equation*} E = \{(x,y) \in \mathbb{R}^m \times \mathbb{R}^n : x_1 y_1 = 1,\ x_2 = \cdots = x_m = 0,\ y_2 = \cdots = y_n = 0\}, \end{equation*}

which is an intersection of finitely many closed sets and hence closed. Indeed \(\{(x,y) : x_1 y_1 = 1\}\) is closed in \(\mathbb{R}^{m+n}\) by exactly the estimate above, because \(\|(x,y) - (x^0,y^0)\|_\infty < \delta\) forces both \(|x_1 - x_1{}^0| < \delta\) and \(|y_1 - y_1{}^0| < \delta\); each set \(\{(x,y) : x_j = 0\}\) is closed because a point with \(x_j \neq 0\) has the open cube \(B((x,y), |x_j|)\) inside the complement, and similarly for each \(\{(x,y) : y_k = 0\}\); and a finite intersection of closed sets is closed because its complement is a finite union of open sets. Then

\begin{equation*} A = \{x \in \mathbb{R}^m : x_1 \neq 0,\ x_2 = \cdots = x_m = 0\}, \end{equation*}

which is not closed: the origin lies in \(\mathbb{R}^m \setminus A\), yet every open cube \(B(0,\delta)\) contains the point \((\delta/2, 0,\dots,0) \in A\), so \(\mathbb{R}^m \setminus A\) is not open.

Problem (5C.12)

(a) Prove that \(\lim_{n \to \infty} \lambda_n(B_n) = 0\).

(b) Find the value of \(n\) that maximizes \(\lambda_n(B_n)\).

Solution

The key tool is the recursion established in the course of proving 5.44: for every integer \(n > 2\),

\begin{equation*} \lambda_n(B_n) = \frac{2\pi}{n}\,\lambda_{n-2}(B_{n-2}). \end{equation*}

We also use the values supplied by 5.44:

\begin{equation*} \lambda_1(B_1) = 2,\quad \lambda_2(B_2) = \pi,\quad \lambda_3(B_3) = \frac{4\pi}{3},\quad \lambda_4(B_4) = \frac{\pi^2}{2},\quad \lambda_5(B_5) = \frac{8\pi^2}{15},\quad \lambda_6(B_6) = \frac{\pi^3}{6}. \end{equation*}

(For example \(\lambda_5(B_5) = 2^{3}\pi^{2}/(1 \cdot 3 \cdot 5) = 8\pi^2/15\) from the odd case of 5.44, and \(\lambda_6(B_6) = \pi^{3}/3! = \pi^3/6\) from the even case.)

(a) Since \(\pi < 3.5\) we have \(2\pi < 7\), so

\begin{equation*} r := \frac{2\pi}{7} < 1 . \end{equation*}

If \(n \ge 5\), then \(n + 2 \ge 7\), so the recursion gives

\begin{equation*} \lambda_{n+2}(B_{n+2}) = \frac{2\pi}{n+2}\,\lambda_n(B_n) \le \frac{2\pi}{7}\,\lambda_n(B_n) = r\,\lambda_n(B_n). \end{equation*}

Iterating this inequality starting from \(n = 5\) and from \(n = 6\) yields, for every integer \(j \ge 0\),

\begin{equation*} \lambda_{5+2j}(B_{5+2j}) \le r^{\,j}\,\lambda_5(B_5) \qquad\text{and}\qquad \lambda_{6+2j}(B_{6+2j}) \le r^{\,j}\,\lambda_6(B_6). \end{equation*}

Let \(M = \max\{\lambda_5(B_5), \lambda_6(B_6)\}\). Every integer \(n \ge 5\) equals \(5 + 2j\) or \(6 + 2j\) for the integer \(j = \lfloor (n-5)/2 \rfloor\), and \(j \ge (n-6)/2\). Hence for all \(n \ge 5\),

\begin{equation*} 0 \le \lambda_n(B_n) \le M\,r^{\,(n-6)/2}. \end{equation*}

Because \(0 < r < 1\), we have \(r^{(n-6)/2} \to 0\) as \(n \to \infty\). The squeeze theorem now gives

\begin{equation*} \lim_{n \to \infty} \lambda_n(B_n) = 0 . \end{equation*}

(b) The answer is \(n = 5\), where \(\lambda_5(B_5) = 8\pi^2/15 \approx 5.26\).

First, \(\lambda_n(B_n)\) is strictly decreasing for \(n \ge 5\) within each parity class: if \(n \ge 5\), then \(n + 2 \ge 7 > 2\pi\), so

\begin{equation*} \lambda_{n+2}(B_{n+2}) = \frac{2\pi}{n+2}\,\lambda_n(B_n) < \lambda_n(B_n), \end{equation*}

using \(\lambda_n(B_n) > 0\). Consequently

\begin{equation*} \lambda_5(B_5) > \lambda_7(B_7) > \lambda_9(B_9) > \cdots \qquad\text{and}\qquad \lambda_6(B_6) > \lambda_8(B_8) > \lambda_{10}(B_{10}) > \cdots . \end{equation*}

So \(\sup_n \lambda_n(B_n) = \max\{\lambda_1(B_1), \lambda_2(B_2), \lambda_3(B_3), \lambda_4(B_4), \lambda_5(B_5), \lambda_6(B_6)\}\), and it suffices to compare those six numbers.

Using \(3 < \pi < 3.2\):

\(\lambda_1(B_1) = 2 < \pi = \lambda_2(B_2)\).

\(\lambda_2(B_2) = \pi < \tfrac{4}{3}\pi = \lambda_3(B_3)\).

\(\lambda_3(B_3) < \lambda_4(B_4)\) because \(\frac{\pi^2/2}{4\pi/3} = \frac{3\pi}{8} > \frac{3 \cdot 3}{8} > 1\).

\(\lambda_4(B_4) < \lambda_5(B_5)\) because \(\frac{8\pi^2/15}{\pi^2/2} = \frac{16}{15} > 1\).

\(\lambda_6(B_6) < \lambda_5(B_5)\) because \(\frac{8\pi^2/15}{\pi^3/6} = \frac{48}{15\pi} = \frac{16}{5\pi}\), and \(5\pi < 16\) since \(\pi < 3.2\).

Thus the largest of the six is \(\lambda_5(B_5)\), and therefore \(\lambda_n(B_n)\) attains its maximum over all positive integers \(n\) uniquely at \(n = 5\), with value \(8\pi^2/15\).

It is worth noting the contrast with the smallest cube containing \(B_n\), which has \(n\)-dimensional Lebesgue measure \(2^n \to \infty\). So the ratio \(\lambda_n(B_n)/2^n\) tends to \(0\) very fast: in high dimensions essentially all of the cube lies outside the inscribed ball.

Problem (5C.13)

For readers familiar with the gamma function \(\Gamma\): Prove that

\begin{equation*} \lambda_n(B_n) = \frac{\pi^{n/2}}{\Gamma\bigl(\frac{n}{2} + 1\bigr)} \end{equation*}

for every positive integer \(n\).

Solution

Recall the two standard properties of the gamma function \(\Gamma(x) = \int_0^\infty t^{x-1}e^{-t}\,dt\) for \(x > 0\) that we need:

\begin{equation*} \Gamma(x+1) = x\,\Gamma(x) \quad \text{for } x > 0, \qquad \Gamma(1) = 1, \qquad \Gamma\bigl(\tfrac12\bigr) = \sqrt{\pi}. \end{equation*}

We argue by induction on \(n\), in steps of \(2\), using the recursion derived in the proof of 5.44: for every integer \(n > 2\),

\begin{equation*} \lambda_n(B_n) = \frac{2\pi}{n}\,\lambda_{n-2}(B_{n-2}). \end{equation*}

Base case \(n = 1\). Here \(\lambda_1(B_1) = 2\), since \(B_1 = (-1,1)\). On the other side,

\begin{equation*} \Gamma\bigl(\tfrac12 + 1\bigr) = \tfrac12\,\Gamma\bigl(\tfrac12\bigr) = \frac{\sqrt{\pi}}{2}, \end{equation*}

so

\begin{equation*} \frac{\pi^{1/2}}{\Gamma\bigl(\frac12 + 1\bigr)} = \frac{\sqrt{\pi}}{\sqrt{\pi}/2} = 2 = \lambda_1(B_1). \end{equation*}

Base case \(n = 2\). Here \(\lambda_2(B_2) = \pi\), and

\begin{equation*} \frac{\pi^{2/2}}{\Gamma\bigl(\frac{2}{2} + 1\bigr)} = \frac{\pi}{\Gamma(2)} = \frac{\pi}{1 \cdot \Gamma(1)} = \pi = \lambda_2(B_2). \end{equation*}

Induction step. Let \(n > 2\) be an integer and suppose the formula holds for \(n - 2\) (which is a positive integer, so the induction hypothesis applies, and every positive integer is reached from \(n = 1\) or \(n = 2\) by repeatedly adding \(2\)). Then, by the recursion and the induction hypothesis,

\begin{equation*} \lambda_n(B_n) = \frac{2\pi}{n} \cdot \frac{\pi^{(n-2)/2}}{\Gamma\bigl(\frac{n-2}{2} + 1\bigr)} = \frac{2\pi \cdot \pi^{(n-2)/2}}{n\,\Gamma\bigl(\frac{n}{2}\bigr)} = \frac{2\,\pi^{n/2}}{n\,\Gamma\bigl(\frac{n}{2}\bigr)} . \end{equation*}

Now apply the functional equation with \(x = n/2 > 0\):

\begin{equation*} \Gamma\Bigl(\frac{n}{2} + 1\Bigr) = \frac{n}{2}\,\Gamma\Bigl(\frac{n}{2}\Bigr), \qquad\text{so}\qquad n\,\Gamma\Bigl(\frac{n}{2}\Bigr) = 2\,\Gamma\Bigl(\frac{n}{2} + 1\Bigr). \end{equation*}

Substituting gives

\begin{equation*} \lambda_n(B_n) = \frac{2\,\pi^{n/2}}{2\,\Gamma\bigl(\frac{n}{2} + 1\bigr)} = \frac{\pi^{n/2}}{\Gamma\bigl(\frac{n}{2} + 1\bigr)}, \end{equation*}

completing the induction.

Hence \(\lambda_n(B_n) = \pi^{n/2}/\Gamma(\frac{n}{2}+1)\) for every positive integer \(n\).

As a consistency check, this single formula reproduces both cases of 5.44. If \(n\) is even, then \(\frac{n}{2}\) is a positive integer and repeated use of \(\Gamma(x+1) = x\Gamma(x)\) together with \(\Gamma(1) = 1\) gives \(\Gamma(\frac{n}{2}+1) = (\frac{n}{2})!\), so \(\lambda_n(B_n) = \pi^{n/2}/(n/2)!\). If \(n\) is odd, write \(n = 2k+1\); then

\begin{equation*} \Gamma\Bigl(\frac{n}{2}+1\Bigr) = \frac{2k+1}{2}\cdot\frac{2k-1}{2}\cdots\frac{3}{2}\cdot\frac{1}{2}\cdot\Gamma\Bigl(\frac12\Bigr) = \frac{1 \cdot 3 \cdot 5 \cdots n}{2^{(n+1)/2}}\sqrt{\pi}, \end{equation*}

so

\begin{equation*} \lambda_n(B_n) = \frac{\pi^{n/2}\,2^{(n+1)/2}}{1 \cdot 3 \cdot 5 \cdots n \cdot \sqrt{\pi}} = \frac{2^{(n+1)/2}\pi^{(n-1)/2}}{1 \cdot 3 \cdot 5 \cdots n}, \end{equation*}

matching the odd case of 5.44.

Problem (5C.14)

Define \(f : \mathbb{R}^2 \to \mathbb{R}\) by

\begin{equation*} f(x,y) = \begin{cases} \dfrac{xy(x^2 - y^2)}{x^2 + y^2} & \text{if } (x,y) \neq (0,0),\\[2ex] 0 & \text{if } (x,y) = (0,0). \end{cases} \end{equation*}

(a) Prove that \(D_1(D_2 f)\) and \(D_2(D_1 f)\) exist everywhere on \(\mathbb{R}^2\).

(b) Show that \(\bigl(D_1(D_2 f)\bigr)(0,0) \neq \bigl(D_2(D_1 f)\bigr)(0,0)\).

(c) Explain why (b) does not violate 5.48.

Solution

Write \(G = \mathbb{R}^2 \setminus \{(0,0)\}\), an open subset of \(\mathbb{R}^2\). On \(G\) we have

\begin{equation*} f(x,y) = \frac{x^3 y - x y^3}{x^2 + y^2}, \end{equation*}

a quotient of polynomials whose denominator does not vanish. Hence \(f\) is infinitely differentiable on \(G\), and in particular all the partial derivatives appearing below exist and are continuous at every point of \(G\).

Step 1: the first-order partial derivatives on all of \(\mathbb{R}^2\).

On \(G\), the quotient rule gives

\begin{equation*} (D_1 f)(x,y) = \frac{(3x^2y - y^3)(x^2+y^2) - (x^3y - xy^3)\cdot 2x}{(x^2+y^2)^2} = \frac{x^4 y + 4x^2 y^3 - y^5}{(x^2+y^2)^2} \end{equation*}

and

\begin{equation*} (D_2 f)(x,y) = \frac{(x^3 - 3xy^2)(x^2+y^2) - (x^3y - xy^3)\cdot 2y}{(x^2+y^2)^2} = \frac{x^5 - 4x^3 y^2 - x y^4}{(x^2+y^2)^2}. \end{equation*}

(Expanding the numerators: \((3x^2y-y^3)(x^2+y^2) = 3x^4y + 2x^2y^3 - y^5\) and \(2x(x^3y-xy^3) = 2x^4y - 2x^2y^3\), whose difference is \(x^4y + 4x^2y^3 - y^5\); similarly \((x^3-3xy^2)(x^2+y^2) = x^5 - 2x^3y^2 - 3xy^4\) and \(2y(x^3y-xy^3) = 2x^3y^2 - 2xy^4\), whose difference is \(x^5 - 4x^3y^2 - xy^4\).)

At the origin, note that \(f(t,0) = 0\) and \(f(0,t) = 0\) for every \(t \in \mathbb{R}\) (this is immediate from the formula when \(t \neq 0\), and from the definition when \(t = 0\)). Hence

\begin{equation*} (D_1 f)(0,0) = \lim_{t \to 0} \frac{f(t,0) - f(0,0)}{t} = 0, \qquad (D_2 f)(0,0) = \lim_{t \to 0} \frac{f(0,t) - f(0,0)}{t} = 0 . \end{equation*}

So \(D_1 f\) and \(D_2 f\) exist at every point of \(\mathbb{R}^2\).

Step 2: the restrictions of \(D_1 f\) and \(D_2 f\) to the axes.

For \(y \neq 0\) the formula for \(D_1 f\) gives

\begin{equation*} (D_1 f)(0,y) = \frac{0 + 0 - y^5}{(y^2)^2} = -y, \end{equation*}

and \((D_1 f)(0,0) = 0 = -0\). Hence

\begin{equation*} (D_1 f)(0,y) = -y \quad \text{for all } y \in \mathbb{R}. \end{equation*}

Likewise, for \(x \neq 0\),

\begin{equation*} (D_2 f)(x,0) = \frac{x^5 - 0 - 0}{(x^2)^2} = x, \end{equation*}

and \((D_2 f)(0,0) = 0\), so

\begin{equation*} (D_2 f)(x,0) = x \quad \text{for all } x \in \mathbb{R}. \end{equation*}

(a) Existence of the two mixed partial derivatives everywhere.

On \(G\) this is immediate: \(D_1 f\) and \(D_2 f\) agree there with the rational functions displayed in Step 1, which are infinitely differentiable on \(G\); hence \(D_1(D_2 f)\) and \(D_2(D_1 f)\) exist at every point of \(G\). (Explicitly, differentiating the formula for \(D_2 f\) with respect to \(x\) and simplifying gives, for \((x,y) \in G\),

\begin{equation*} \bigl(D_1(D_2 f)\bigr)(x,y) = \frac{(5x^4 - 12x^2y^2 - y^4)(x^2+y^2) - 4x^2(x^4 - 4x^2y^2 - y^4)}{(x^2+y^2)^3} = \frac{x^6 + 9x^4 y^2 - 9x^2 y^4 - y^6}{(x^2+y^2)^3}, \end{equation*}

and the same expression is obtained for \(D_2(D_1 f)\) on \(G\), as it must be by 5.48 applied on the open set \(G\), where \(f\) is smooth.)

At the origin, by definition

\begin{equation*} \bigl(D_1(D_2 f)\bigr)(0,0) = \lim_{t \to 0} \frac{(D_2 f)(t,0) - (D_2 f)(0,0)}{t}, \end{equation*}

which makes sense because \(D_2 f\) exists everywhere by Step 1. By Step 2, \((D_2 f)(t,0) = t\), so this limit equals \(\lim_{t \to 0} t/t = 1\); in particular it exists. Similarly

\begin{equation*} \bigl(D_2(D_1 f)\bigr)(0,0) = \lim_{t \to 0} \frac{(D_1 f)(0,t) - (D_1 f)(0,0)}{t} = \lim_{t \to 0} \frac{-t}{t} = -1, \end{equation*}

which also exists. Thus \(D_1(D_2 f)\) and \(D_2(D_1 f)\) exist at every point of \(\mathbb{R}^2\).

(b) From the computation just performed,

\begin{equation*} \bigl(D_1(D_2 f)\bigr)(0,0) = 1 \neq -1 = \bigl(D_2(D_1 f)\bigr)(0,0). \end{equation*}

(c) Result 5.48 requires more than the mere existence of \(D_1 f\), \(D_2 f\), \(D_1(D_2 f)\), and \(D_2(D_1 f)\) on the open set \(G\) in its statement: it requires all four of these functions to be continuous on that open set. For our \(f\), at least one of \(D_1(D_2 f)\), \(D_2(D_1 f)\) fails to be continuous at the origin, so 5.48 simply does not apply on any open set containing \((0,0)\).

Here is a proof of that failure that uses no computation beyond part (b). On the open set \(G\) the function \(f\) is infinitely differentiable, so all four functions \(D_1 f\), \(D_2 f\), \(D_1(D_2 f)\), \(D_2(D_1 f)\) are continuous on \(G\), and 5.48 applies there to give

\begin{equation*} D_1(D_2 f) = D_2(D_1 f) \quad \text{on } G . \end{equation*}

If both mixed partials were also continuous at \((0,0)\), then taking the limit as \((x,y) \to (0,0)\) within \(G\) of the two equal functions would give \(\bigl(D_1(D_2 f)\bigr)(0,0) = \bigl(D_2(D_1 f)\bigr)(0,0)\), contradicting (b). Hence continuity of the mixed partials at the origin fails, and the hypotheses of 5.48 are not satisfied on any open set containing the origin.

One can also see the discontinuity explicitly from the formula in (a). For \(x \neq 0\),

\begin{equation*} \bigl(D_1(D_2 f)\bigr)(x,0) = \frac{x^6}{x^6} = 1, \qquad\text{while for } y \neq 0, \quad \bigl(D_1(D_2 f)\bigr)(0,y) = \frac{-y^6}{y^6} = -1 . \end{equation*}

So \(D_1(D_2 f)\) takes the value \(1\) at points arbitrarily close to the origin along the horizontal axis and the value \(-1\) at points arbitrarily close to the origin along the vertical axis; since \(\bigl(D_1(D_2 f)\bigr)(0,0) = 1\), no limit at the origin exists and \(D_1(D_2 f)\) is discontinuous there. The same computation shows \(D_2(D_1 f)\), which equals the same rational function on \(G\) but has value \(-1\) at the origin, is discontinuous there as well.

This example is exactly the one referred to in the marginal comment preceding 5.48: the continuity hypotheses in 5.48 can be weakened somewhat, but they cannot be eliminated.

Banach Spaces

Exercises 6A

Problem (6A.1)

Verify that each of the claimed metrics in Example 6.2 is indeed a metric.

Solution

Recall (6.1) that a function \(d \colon V \times V \to [0,\infty)\) is a metric provided that \(d(f,f) = 0\) for all \(f\), that \(d(f,g) = 0\) implies \(f = g\), that \(d(f,g) = d(g,f)\), and that \(d(f,h) \le d(f,g) + d(g,h)\). We check these four conditions for each of the five examples in 6.2. In every case the first three conditions are nearly immediate; the content is in checking that \(d\) really takes values in \([0,\infty)\) (finiteness) and in the triangle inequality.

The discrete metric. Suppose \(V\) is a nonempty set and \(d(f,g) = 1\) if \(f \ne g\) and \(d(f,g) = 0\) if \(f = g\). Then \(d\) takes values in \(\{0,1\} \subseteq [0,\infty)\). By definition \(d(f,f) = 0\), and if \(d(f,g) = 0\) then we cannot have \(f \ne g\) (which would force \(d(f,g) = 1\)), so \(f = g\). Because the relation \(f \ne g\) is symmetric in \(f\) and \(g\), we have \(d(f,g) = d(g,f)\). For the triangle inequality, suppose \(f, g, h \in V\). If \(f = h\) then \(d(f,h) = 0 \le d(f,g) + d(g,h)\) because \(d\) is nonnegative. If \(f \ne h\), then \(g\) cannot equal both \(f\) and \(h\); hence at least one of \(d(f,g)\), \(d(g,h)\) equals \(1\), and thus

\begin{equation*} d(f,h) = 1 \le d(f,g) + d(g,h). \end{equation*}

The absolute-value metric on \(\mathbb{R}\). Here \(d(x,y) = |x-y| \in [0,\infty)\). We have \(d(x,x) = |x - x| = 0\); if \(|x - y| = 0\) then \(x - y = 0\), so \(x = y\); and \(|x-y| = |-(x-y)| = |y-x|\). For \(x,y,z \in \mathbb{R}\), the triangle inequality for absolute value gives

\begin{equation*} d(x,z) = |x - z| = |(x-y) + (y-z)| \le |x-y| + |y-z| = d(x,y) + d(y,z). \end{equation*}

The max metric on \(\mathbb{R}^n\). For \(x = (x_1,\dots,x_n)\) and \(y = (y_1,\dots,y_n)\) in \(\mathbb{R}^n\), set \(d(x,y) = \max\{|x_1-y_1|,\dots,|x_n-y_n|\}\). This is a maximum of finitely many nonnegative real numbers, hence is a well-defined element of \([0,\infty)\). Clearly \(d(x,x) = \max\{0,\dots,0\} = 0\). If \(d(x,y) = 0\) then \(|x_j - y_j| \le 0\) for each \(j \in \{1,\dots,n\}\), so \(x_j = y_j\) for each \(j\) and thus \(x = y\). Symmetry holds because \(|x_j - y_j| = |y_j - x_j|\) for each \(j\). For the triangle inequality, let \(x, y, z \in \mathbb{R}^n\) and let \(j \in \{1,\dots,n\}\). Then

\begin{equation*} |x_j - z_j| \le |x_j - y_j| + |y_j - z_j| \le d(x,y) + d(y,z), \end{equation*}

where the second inequality holds because each summand is at most the corresponding maximum. The right side does not depend on \(j\), so taking the maximum over \(j \in \{1,\dots,n\}\) on the left gives \(d(x,z) \le d(x,y) + d(y,z)\).

The sup metric on \(C([0,1])\). Here \(C([0,1])\) is the set of continuous real-valued functions on \([0,1]\) and \(d(f,g) = \sup\{|f(t) - g(t)| : t \in [0,1]\}\). First we check that \(d(f,g) \in [0,\infty)\). If \(f, g \in C([0,1])\) then \(|f - g|\) is a continuous real-valued function on the closed bounded interval \([0,1]\) and hence is bounded; thus the set \(\{|f(t)-g(t)| : t \in [0,1]\}\) is a nonempty bounded set of nonnegative numbers and its supremum is a nonnegative real number. Next, \(d(f,f) = \sup\{0\} = 0\). If \(d(f,g) = 0\) then \(|f(t) - g(t)| \le 0\) for every \(t \in [0,1]\), so \(f(t) = g(t)\) for every \(t \in [0,1]\), which means \(f = g\). Symmetry holds because \(|f(t) - g(t)| = |g(t) - f(t)|\) for each \(t\). For the triangle inequality, suppose \(f, g, h \in C([0,1])\) and \(t \in [0,1]\). Then

\begin{equation*} |f(t) - h(t)| \le |f(t) - g(t)| + |g(t) - h(t)| \le d(f,g) + d(g,h). \end{equation*}

Thus \(d(f,g) + d(g,h)\) is an upper bound for \(\{|f(t)-h(t)| : t \in [0,1]\}\), and therefore \(d(f,h) \le d(f,g) + d(g,h)\).

The \(\ell^1\) metric. Here \(\ell^1\) is the set of sequences \((a_1,a_2,\dots)\) of real numbers with \(\sum_{k=1}^{\infty} |a_k| < \infty\), and

\begin{equation*} d\bigl((a_1,a_2,\dots),(b_1,b_2,\dots)\bigr) = \sum_{k=1}^{\infty} |a_k - b_k|. \end{equation*}

To see that \(d\) takes values in \([0,\infty)\), note that all terms are nonnegative and that \(|a_k - b_k| \le |a_k| + |b_k|\) for each \(k\), so

\begin{equation*} \sum_{k=1}^{\infty} |a_k - b_k| \le \sum_{k=1}^{\infty} |a_k| + \sum_{k=1}^{\infty} |b_k| < \infty \end{equation*}

whenever \((a_k)\) and \((b_k)\) are in \(\ell^1\). Clearly \(d(a,a) = 0\). If \(d(a,b) = 0\) then, because every term of the series is nonnegative, \(|a_k - b_k| = 0\) for every \(k\); hence \(a_k = b_k\) for every \(k\) and the two sequences are equal. Symmetry is immediate from \(|a_k - b_k| = |b_k - a_k|\). Finally, suppose \(a, b, c \in \ell^1\). For every \(k\) we have \(|a_k - c_k| \le |a_k - b_k| + |b_k - c_k|\), so for every \(n \in \mathbb{Z}^+\) the partial sums satisfy

\begin{equation*} \sum_{k=1}^{n} |a_k - c_k| \le \sum_{k=1}^{n} |a_k - b_k| + \sum_{k=1}^{n} |b_k - c_k| \le d(a,b) + d(b,c), \end{equation*}

where the last inequality holds because each of the two series has nonnegative terms and hence has partial sums bounded by its sum. Letting \(n \to \infty\) gives

\begin{equation*} d(a,c) = \sum_{k=1}^{\infty} |a_k - c_k| \le d(a,b) + d(b,c). \end{equation*}

Thus each of the five functions in 6.2 is a metric on the indicated set.

Problem (6A.2)

Prove that every finite subset of a metric space is closed.

Solution

Suppose \((V,d)\) is a metric space and \(E \subseteq V\) is finite. By 6.6 we must show that \(V \setminus E\) is an open subset of \(V\).

First suppose \(E = \emptyset\). Then \(V \setminus E = V\), and \(V\) is open because \(B(f,1) \subseteq V\) for every \(f \in V\). Hence \(\emptyset\) is closed.

Now suppose \(E \ne \emptyset\), say \(E = \{f_1,\dots,f_n\}\) with \(n \in \mathbb{Z}^+\). Suppose \(g \in V \setminus E\). Then \(g \ne f_j\) for each \(j \in \{1,\dots,n\}\). The second bullet point in 6.1 (in contrapositive form) shows that \(d(g,f_j) \ne 0\) for each such \(j\); because \(d\) takes values in \([0,\infty)\), this means \(d(g,f_j) > 0\) for each \(j \in \{1,\dots,n\}\). Let

\begin{equation*} r = \min\{d(g,f_1),\dots,d(g,f_n)\}. \end{equation*}

As the minimum of a finite nonempty collection of positive numbers, \(r > 0\).

We claim \(B(g,r) \subseteq V \setminus E\). Indeed, if some \(h \in B(g,r)\) were an element of \(E\), then \(h = f_j\) for some \(j \in \{1,\dots,n\}\), and hence

\begin{equation*} d(g,f_j) = d(g,h) < r \le d(g,f_j), \end{equation*}

a contradiction. Thus \(B(g,r) \cap E = \emptyset\), which is to say \(B(g,r) \subseteq V \setminus E\).

Because every \(g \in V \setminus E\) is the center of an open ball contained in \(V \setminus E\), the set \(V \setminus E\) is open (6.4). Hence \(E\) is closed.

Problem (6A.3)

Prove that every closed ball in a metric space is closed.

Solution

Suppose \((V,d)\) is a metric space, \(f \in V\), and \(r > 0\). Recall from 6.3 that

\begin{equation*} \overline{B}(f,r) = \{g \in V : d(f,g) \le r\}. \end{equation*}

By 6.6 we must show that \(V \setminus \overline{B}(f,r)\) is open.

Suppose \(h \in V \setminus \overline{B}(f,r)\), which means \(d(f,h) > r\). Let

\begin{equation*} s = d(f,h) - r, \end{equation*}

so \(s > 0\). We claim that \(B(h,s) \subseteq V \setminus \overline{B}(f,r)\).

To see this, suppose \(g \in B(h,s)\), so \(d(h,g) < s\). The triangle inequality (the last bullet point of 6.1), together with the symmetry of \(d\), gives

\begin{equation*} d(f,h) \le d(f,g) + d(g,h) < d(f,g) + s. \end{equation*}

Hence

\begin{equation*} d(f,g) > d(f,h) - s = d(f,h) - \bigl(d(f,h) - r\bigr) = r, \end{equation*}

so \(g \notin \overline{B}(f,r)\). This proves the claim.

Thus every element of \(V \setminus \overline{B}(f,r)\) is the center of an open ball contained in \(V \setminus \overline{B}(f,r)\), so \(V \setminus \overline{B}(f,r)\) is open by 6.4. Therefore \(\overline{B}(f,r)\) is a closed subset of \(V\).

Problem (6A.4)

Suppose \(V\) is a metric space.

(a) Prove that the union of each collection of open subsets of \(V\) is an open subset of \(V\).

(b) Prove that the intersection of each finite collection of open subsets of \(V\) is an open subset of \(V\).

Solution

Throughout, \(d\) denotes the metric on \(V\), and we use the definition of open set given by 6.4: a subset \(G\) of \(V\) is open if for every \(f \in G\) there exists \(r > 0\) with \(B(f,r) \subseteq G\).

(a) Suppose \(\mathcal{A}\) is a collection of open subsets of \(V\) and let

\begin{equation*} G = \bigcup_{A \in \mathcal{A}} A . \end{equation*}

Suppose \(f \in G\). Then there exists \(A \in \mathcal{A}\) with \(f \in A\). Because \(A\) is open, there exists \(r > 0\) such that \(B(f,r) \subseteq A\). Since \(A \subseteq G\), we conclude that \(B(f,r) \subseteq G\). As \(f\) was an arbitrary element of \(G\), this shows that \(G\) is open.

(Note that if \(\mathcal{A} = \emptyset\) then \(G = \emptyset\) and the condition in 6.4 holds vacuously, so the empty set is open; the argument above covers this case automatically.)

(b) Suppose \(G_1,\dots,G_n\) are open subsets of \(V\), where \(n \in \mathbb{Z}^+\), and let

\begin{equation*} G = G_1 \cap \cdots \cap G_n . \end{equation*}

Suppose \(f \in G\). Then \(f \in G_j\) for each \(j \in \{1,\dots,n\}\), so for each such \(j\) the openness of \(G_j\) provides \(r_j > 0\) with \(B(f,r_j) \subseteq G_j\). Let

\begin{equation*} r = \min\{r_1,\dots,r_n\}, \end{equation*}

which is positive because it is the minimum of finitely many positive numbers. If \(g \in B(f,r)\) and \(j \in \{1,\dots,n\}\), then \(d(f,g) < r \le r_j\), so \(g \in B(f,r_j) \subseteq G_j\). Hence \(B(f,r) \subseteq G_j\) for every \(j\), and therefore \(B(f,r) \subseteq G\). As \(f\) was an arbitrary element of \(G\), this shows that \(G\) is open.

(If one adopts the convention that the intersection of the empty collection of subsets of \(V\) equals \(V\), then that case is also fine: \(V\) is open because \(B(f,1) \subseteq V\) for every \(f \in V\).)

The hypothesis that the collection in (b) is finite cannot be dropped: in \(\mathbb{R}\) with the metric \(d(x,y) = |x-y|\), each set \(\bigl(-\frac{1}{k},\frac{1}{k}\bigr)\) is open, but \(\bigcap_{k=1}^{\infty}\bigl(-\frac{1}{k},\frac{1}{k}\bigr) = \{0\}\), which is not open.

Problem (6A.5)

Suppose \(V\) is a metric space.

(a) Prove that the intersection of each collection of closed subsets of \(V\) is a closed subset of \(V\).

(b) Prove that the union of each finite collection of closed subsets of \(V\) is a closed subset of \(V\).

Solution

By 6.6, a subset of \(V\) is closed exactly when its complement in \(V\) is open. Thus both parts follow from Exercise 6A.4 by taking complements and using De Morgan’s laws.

(a) Suppose \(\mathcal{F}\) is a nonempty collection of closed subsets of \(V\) and let

\begin{equation*} F = \bigcap_{A \in \mathcal{F}} A . \end{equation*}

De Morgan’s laws give

\begin{equation*} V \setminus F = V \setminus \bigcap_{A \in \mathcal{F}} A = \bigcup_{A \in \mathcal{F}} (V \setminus A). \end{equation*}

Indeed, \(f \in V \setminus F\) if and only if \(f \in V\) and \(f\) fails to belong to at least one \(A \in \mathcal{F}\), which happens if and only if \(f \in V \setminus A\) for some \(A \in \mathcal{F}\).

Because each \(A \in \mathcal{F}\) is closed, each set \(V \setminus A\) is open. Hence \(V \setminus F\) is a union of a collection of open subsets of \(V\), so \(V \setminus F\) is open by Exercise 6A.4(a). Therefore \(F\) is closed.

(If \(\mathcal{F} = \emptyset\) and one interprets the intersection as \(V\), the conclusion still holds: \(V \setminus V = \emptyset\) is open, so \(V\) is closed.)

(b) Suppose \(F_1,\dots,F_n\) are closed subsets of \(V\), where \(n \in \mathbb{Z}^+\), and let

\begin{equation*} F = F_1 \cup \cdots \cup F_n . \end{equation*}

De Morgan’s laws give

\begin{equation*} V \setminus F = V \setminus \bigcup_{j=1}^{n} F_j = \bigcap_{j=1}^{n} (V \setminus F_j), \end{equation*}

since an element of \(V\) lies outside \(F\) if and only if it lies outside every \(F_j\).

Because each \(F_j\) is closed, each set \(V \setminus F_j\) is open. Hence \(V \setminus F\) is the intersection of a finite collection of open subsets of \(V\), so \(V \setminus F\) is open by Exercise 6A.4(b). Therefore \(F\) is closed.

(The union of the empty collection is \(\emptyset\), whose complement \(V\) is open, so \(\emptyset\) is closed as well.)

Again finiteness in (b) is essential: in \(\mathbb{R}\) with \(d(x,y) = |x-y|\), each singleton \(\{1/k\}\) is closed by Exercise 6A.2, but \(\bigcup_{k=1}^{\infty}\{1/k\}\) is not closed, since \(0\) lies in its closure (it is the limit of the sequence \(1, \frac12, \frac13, \dots\)) but not in the set; apply 6.9(e).

Problem (6A.6)

(a) Prove that if \(V\) is a metric space, \(f \in V\), and \(r > 0\), then \(\overline{B(f,r)} \subseteq \overline{B}(f,r)\).

(b) Give an example of a metric space \(V\), \(f \in V\), and \(r > 0\) such that \(\overline{B(f,r)} \ne \overline{B}(f,r)\).

Solution

Here \(B(f,r) = \{g \in V : d(f,g) < r\}\) is the open ball and \(\overline{B}(f,r) = \{g \in V : d(f,g) \le r\}\) is the closed ball (6.3), while \(\overline{B(f,r)}\) denotes the closure of the open ball in the sense of 6.7.

(a) Suppose \(g \in \overline{B(f,r)}\). By 6.7 this means

\begin{equation*} B(g,\varepsilon) \cap B(f,r) \ne \emptyset \quad \text{for every } \varepsilon > 0 . \end{equation*}

Let \(\varepsilon > 0\) and choose \(h \in B(g,\varepsilon) \cap B(f,r)\). Then \(d(g,h) < \varepsilon\) and \(d(f,h) < r\), so the triangle inequality and the symmetry of \(d\) (6.1) give

\begin{equation*} d(f,g) \le d(f,h) + d(h,g) < r + \varepsilon . \end{equation*}

Thus \(d(f,g) < r + \varepsilon\) for every \(\varepsilon > 0\), which forces \(d(f,g) \le r\). Hence \(g \in \overline{B}(f,r)\), completing the proof that \(\overline{B(f,r)} \subseteq \overline{B}(f,r)\).

Alternatively, one can argue as follows: \(\overline{B}(f,r)\) is a closed subset of \(V\) by Exercise 6A.3, and it contains \(B(f,r)\) because \(d(f,g) < r\) implies \(d(f,g) \le r\); since by 6.9(b) the closure \(\overline{B(f,r)}\) is the intersection of all closed subsets of \(V\) containing \(B(f,r)\), we get \(\overline{B(f,r)} \subseteq \overline{B}(f,r)\).

(b) Let \(V = \{0,1\}\) with the discrete metric from the first bullet point of 6.2, so \(d(x,y) = 1\) if \(x \ne y\) and \(d(x,y) = 0\) if \(x = y\). Take \(f = 0\) and \(r = 1\).

Then

\begin{equation*} B(0,1) = \{g \in V : d(0,g) < 1\} = \{0\}, \qquad \overline{B}(0,1) = \{g \in V : d(0,g) \le 1\} = \{0,1\} = V, \end{equation*}

because \(d(0,1) = 1\), which is not less than \(1\) but is at most \(1\).

Now \(\overline{B(0,1)} = \overline{\{0\}} = \{0\}\). To see this directly from 6.7, note that \(0 \in \overline{\{0\}}\) trivially, while \(1 \notin \overline{\{0\}}\) because \(B(1,\tfrac12) = \{1\}\) (as \(d(1,0) = 1 > \tfrac12\)), so \(B(1,\tfrac12) \cap \{0\} = \emptyset\). (One can also just invoke Exercise 6A.2 and 6.9(d): the finite set \(\{0\}\) is closed, hence equals its own closure.)

Therefore

\begin{equation*} \overline{B(0,1)} = \{0\} \ne \{0,1\} = \overline{B}(0,1). \end{equation*}

The same phenomenon occurs, for example, in \(V = \mathbb{Z}\) with the metric \(d(m,n) = |m-n|\) inherited from \(\mathbb{R}\): there \(B(0,1) = \{0\}\), whose closure is \(\{0\}\), while \(\overline{B}(0,1) = \{-1,0,1\}\).

Problem (6A.7)

Show that each sequence in a metric space has at most one limit.

Solution

Suppose \((V,d)\) is a metric space, \(f_1, f_2, \dots\) is a sequence in \(V\), and \(f, g \in V\) satisfy

\begin{equation*} \lim_{k \to \infty} f_k = f \quad \text{and} \quad \lim_{k \to \infty} f_k = g . \end{equation*}

We show that \(f = g\); this proves that a sequence cannot have two different limits, which is exactly the assertion that it has at most one limit. (Of course a sequence may have no limit at all.)

Let \(\varepsilon > 0\). By 6.8, \(\lim_{k \to \infty} d(f_k,f) = 0\), so there exists \(n_1 \in \mathbb{Z}^+\) such that

\begin{equation*} d(f_k,f) < \tfrac{\varepsilon}{2} \quad \text{for all integers } k \ge n_1 . \end{equation*}

Similarly there exists \(n_2 \in \mathbb{Z}^+\) such that

\begin{equation*} d(f_k,g) < \tfrac{\varepsilon}{2} \quad \text{for all integers } k \ge n_2 . \end{equation*}

Let \(n = \max\{n_1,n_2\}\). Then both inequalities hold for \(k = n\), and hence the triangle inequality together with the symmetry of \(d\) (6.1) gives

\begin{equation*} d(f,g) \le d(f,f_n) + d(f_n,g) < \tfrac{\varepsilon}{2} + \tfrac{\varepsilon}{2} = \varepsilon . \end{equation*}

Thus \(0 \le d(f,g) < \varepsilon\) for every \(\varepsilon > 0\). If we had \(d(f,g) > 0\), then taking \(\varepsilon = d(f,g)\) would yield the contradiction \(d(f,g) < d(f,g)\). Hence \(d(f,g) = 0\), and the second bullet point of 6.1 implies \(f = g\).

Therefore each sequence in a metric space has at most one limit, which justifies the notation \(\lim_{k \to \infty} f_k\) for the limit of a convergent sequence.

Problem (6A.8)

Prove 6.9.

[6.9 states: Suppose \(V\) is a metric space and \(E \subseteq V\). Then

(a) \(\overline{E} = \{g \in V : \text{there exist } f_1, f_2, \ldots \text{ in } E \text{ such that } \lim_{k \to \infty} f_k = g\}\);

(b) \(\overline{E}\) is the intersection of all closed subsets of \(V\) that contain \(E\);

(c) \(\overline{E}\) is a closed subset of \(V\);

(d) \(E\) is closed if and only if \(E = \overline{E}\);

(e) \(E\) is closed if and only if \(E\) contains the limit of every convergent sequence of elements of \(E\).]

Solution

Throughout, \(d\) denotes the metric on \(V\). Recall the definition 6.7 of the closure:

\begin{equation*} \overline{E} = \{g \in V : B(g,\varepsilon) \cap E \neq \emptyset \text{ for every } \varepsilon > 0\}. \end{equation*}

Preliminary observation: \(E \subseteq \overline{E}\). Indeed, if \(f \in E\) and \(\varepsilon > 0\), then \(d(f,f) = 0 < \varepsilon\), so \(f \in B(f,\varepsilon) \cap E\), showing that this intersection is nonempty; hence \(f \in \overline{E}\).

Proof of (a). Let \(A\) denote the set on the right side of (a).

Suppose \(g \in \overline{E}\). For each \(k \in \mathbb{Z}^+\), the set \(B(g, 1/k) \cap E\) is nonempty; choose \(f_k\) in it. Then \(f_k \in E\) and \(d(f_k, g) < 1/k\) for each \(k\), so \(\lim_{k \to \infty} d(f_k, g) = 0\), which by 6.8 means \(\lim_{k \to \infty} f_k = g\). Thus \(g \in A\).

Conversely, suppose \(g \in A\), so there exist \(f_1, f_2, \ldots\) in \(E\) with \(\lim_{k \to \infty} f_k = g\). Let \(\varepsilon > 0\). Because \(\lim_{k \to \infty} d(f_k, g) = 0\), there exists \(k \in \mathbb{Z}^+\) with \(d(f_k, g) < \varepsilon\). Then \(f_k \in B(g,\varepsilon) \cap E\), so this intersection is nonempty. As \(\varepsilon > 0\) was arbitrary, \(g \in \overline{E}\).

Hence \(\overline{E} = A\).

Proof of (c). We show that \(V \setminus \overline{E}\) is open; then \(\overline{E}\) is closed by 6.6. Suppose \(g \in V \setminus \overline{E}\). By the definition of closure, there exists \(\varepsilon > 0\) such that \(B(g,\varepsilon) \cap E = \emptyset\). We claim \(B(g,\varepsilon) \subseteq V \setminus \overline{E}\). To see this, let \(h \in B(g,\varepsilon)\) and set \(\delta = \varepsilon - d(g,h) > 0\). If \(u \in B(h,\delta)\), then

\begin{equation*} d(g,u) \leq d(g,h) + d(h,u) < d(g,h) + \delta = \varepsilon, \end{equation*}

so \(B(h,\delta) \subseteq B(g,\varepsilon)\) and therefore \(B(h,\delta) \cap E = \emptyset\). Thus \(h \notin \overline{E}\), proving the claim. Hence every element of \(V \setminus \overline{E}\) is the center of an open ball contained in \(V \setminus \overline{E}\), so \(V \setminus \overline{E}\) is open and \(\overline{E}\) is closed.

Proof of (b). First we record the following statement, labelled \((\ast)\): if \(F\) is a closed subset of \(V\) with \(E \subseteq F\), then \(\overline{E} \subseteq F\).

To prove \((\ast)\), suppose \(g \in V \setminus F\). Because \(V \setminus F\) is open, there exists \(\varepsilon > 0\) with \(B(g,\varepsilon) \subseteq V \setminus F \subseteq V \setminus E\), and hence \(B(g,\varepsilon) \cap E = \emptyset\), which gives \(g \notin \overline{E}\). Taking complements yields \(\overline{E} \subseteq F\).

Now let \(C\) be the intersection of all closed subsets of \(V\) that contain \(E\). By \((\ast)\), the set \(\overline{E}\) is contained in each such closed set, so \(\overline{E} \subseteq C\). On the other hand, by the preliminary observation and (c), the set \(\overline{E}\) is itself a closed subset of \(V\) containing \(E\); hence \(\overline{E}\) is one of the sets being intersected, and therefore \(C \subseteq \overline{E}\). Thus \(C = \overline{E}\), as claimed.

Proof of (d). If \(E = \overline{E}\), then \(E\) is closed by (c). Conversely, suppose \(E\) is closed. Then \(E\) is a closed set containing \(E\), so \((\ast)\) gives \(\overline{E} \subseteq E\); combined with the preliminary observation \(E \subseteq \overline{E}\), this gives \(E = \overline{E}\).

Proof of (e). Suppose \(E\) is closed. If \(f_1, f_2, \ldots\) is a sequence in \(E\) converging to some \(g \in V\), then \(g \in \overline{E}\) by (a), and \(\overline{E} = E\) by (d); hence \(g \in E\). Thus \(E\) contains the limit of every convergent sequence of its elements.

Conversely, suppose \(E\) contains the limit of every convergent sequence of elements of \(E\). If \(g \in \overline{E}\), then by (a) there is a sequence of elements of \(E\) converging to \(g\), so \(g \in E\) by hypothesis. Thus \(\overline{E} \subseteq E\), and with the preliminary observation we get \(E = \overline{E}\), so \(E\) is closed by (d).

Problem (6A.9)

Prove that each open subset of a metric space \(V\) is the union of some sequence of closed subsets of \(V\).

Solution

Let \(d\) be the metric on \(V\) and let \(G\) be an open subset of \(V\). For each \(k \in \mathbb{Z}^+\), define

\begin{equation*} F_k = V \setminus \bigcup_{h \in V \setminus G} B\!\left(h, \tfrac{1}{k}\right). \end{equation*}

(If \(G = V\), the union is over the empty index set and is thus \(\emptyset\), so \(F_k = V\) in that case.) Explicitly,

\begin{equation*} F_k = \left\{ f \in V : d(f,h) \geq \tfrac{1}{k} \text{ for every } h \in V \setminus G \right\}. \end{equation*}

Each \(F_k\) is closed. Indeed, each ball \(B(h, 1/k)\) is an open subset of \(V\) by 6.5, and a union of open sets is open: if \(f\) belongs to a union \(\bigcup_{\alpha} G_\alpha\) of open sets, then \(f \in G_\alpha\) for some \(\alpha\), so there exists \(r > 0\) with \(B(f,r) \subseteq G_\alpha \subseteq \bigcup_\alpha G_\alpha\). Hence \(V \setminus F_k\) is open, which by 6.6 means precisely that \(F_k\) is closed.

Next, \(F_k \subseteq G\) for every \(k\). Indeed, if \(f \in V \setminus G\), then \(d(f,f) = 0 < 1/k\), so \(f \in B(f, 1/k)\) and therefore \(f \notin F_k\). Taking contrapositives, \(f \in F_k\) implies \(f \in G\). Consequently

\begin{equation*} \bigcup_{k=1}^{\infty} F_k \subseteq G. \end{equation*}

Finally, \(G \subseteq \bigcup_{k=1}^{\infty} F_k\). Suppose \(f \in G\). Because \(G\) is open, there exists \(r > 0\) with \(B(f,r) \subseteq G\). Choose \(k \in \mathbb{Z}^+\) with \(1/k \leq r\). If we had \(f \in B(h, 1/k)\) for some \(h \in V \setminus G\), then \(d(f,h) < 1/k \leq r\), so \(h \in B(f,r) \subseteq G\), contradicting \(h \notin G\). Hence \(f\) lies in no such ball, which means \(f \in F_k\).

Combining the last two displays, \(G = \bigcup_{k=1}^{\infty} F_k\), exhibiting the open set \(G\) as the union of the sequence \(F_1, F_2, \ldots\) of closed subsets of \(V\).

Problem (6A.10)

Prove or give a counterexample: If \(V\) is a metric space and \(U, W\) are subsets of \(V\), then \(\overline{U} \cup \overline{W} = \overline{U \cup W}\).

Solution

The statement is true; we prove it.

Let \(d\) be the metric on \(V\). Recall from 6.7 that for \(E \subseteq V\),

\begin{equation*} \overline{E} = \{g \in V : B(g,\varepsilon) \cap E \neq \emptyset \text{ for every } \varepsilon > 0\}. \end{equation*}

First note that the closure operation is order preserving: if \(E_1 \subseteq E_2 \subseteq V\) and \(g \in \overline{E_1}\), then for every \(\varepsilon > 0\) we have \(\emptyset \neq B(g,\varepsilon) \cap E_1 \subseteq B(g,\varepsilon) \cap E_2\), so \(g \in \overline{E_2}\). Thus \(E_1 \subseteq E_2\) implies \(\overline{E_1} \subseteq \overline{E_2}\).

Applying this to \(U \subseteq U \cup W\) and to \(W \subseteq U \cup W\) gives \(\overline{U} \subseteq \overline{U \cup W}\) and \(\overline{W} \subseteq \overline{U \cup W}\), hence

\begin{equation*} \overline{U} \cup \overline{W} \subseteq \overline{U \cup W}. \end{equation*}

For the reverse inclusion, suppose \(g \in V\) and \(g \notin \overline{U} \cup \overline{W}\); we show \(g \notin \overline{U \cup W}\). Because \(g \notin \overline{U}\), there exists \(\varepsilon_1 > 0\) with \(B(g,\varepsilon_1) \cap U = \emptyset\). Because \(g \notin \overline{W}\), there exists \(\varepsilon_2 > 0\) with \(B(g,\varepsilon_2) \cap W = \emptyset\). Let \(\varepsilon = \min\{\varepsilon_1, \varepsilon_2\} > 0\). Then \(B(g,\varepsilon) \subseteq B(g,\varepsilon_1)\) and \(B(g,\varepsilon) \subseteq B(g,\varepsilon_2)\), so

\begin{equation*} B(g,\varepsilon) \cap (U \cup W) = \bigl(B(g,\varepsilon) \cap U\bigr) \cup \bigl(B(g,\varepsilon) \cap W\bigr) = \emptyset . \end{equation*}

Hence \(g \notin \overline{U \cup W}\). Taking contrapositives gives \(\overline{U \cup W} \subseteq \overline{U} \cup \overline{W}\).

The two inclusions give \(\overline{U} \cup \overline{W} = \overline{U \cup W}\).

Alternative proof of the inclusion \(\overline{U \cup W} \subseteq \overline{U} \cup \overline{W}\): by 6.9(c) each of \(\overline{U}\) and \(\overline{W}\) is closed, so their union is closed (a union of two closed sets is closed, because the complement of the union is the intersection of the two open complements, which is open). Since \(U \cup W \subseteq \overline{U} \cup \overline{W}\) by the preliminary observation \(E \subseteq \overline{E}\), and since by 6.9(b) the set \(\overline{U \cup W}\) is the intersection of all closed sets containing \(U \cup W\), we conclude \(\overline{U \cup W} \subseteq \overline{U} \cup \overline{W}\).

Note that the corresponding statement for infinite unions is false: in \(\mathbb{R}\) with its usual metric, writing \(\mathbb{Q} = \{r_1, r_2, \ldots\}\), each singleton \(\{r_k\}\) is closed and equals its own closure, but \(\overline{\bigcup_k \{r_k\}} = \overline{\mathbb{Q}} = \mathbb{R} \neq \mathbb{Q} = \bigcup_k \overline{\{r_k\}}\).

Problem (6A.11)

Prove or give a counterexample: If \(V\) is a metric space and \(U, W\) are subsets of \(V\), then \(\overline{U} \cap \overline{W} = \overline{U \cap W}\).

Solution

The statement is false. Here is a counterexample.

Let \(V = \mathbb{R}\) with the usual metric \(d(x,y) = |x - y|\) (see the second bullet point of 6.2), and let

\begin{equation*} U = (0,1), \qquad W = (1,2). \end{equation*}

First, \(U \cap W = \emptyset\), and \(\overline{\emptyset} = \emptyset\): for any \(g \in \mathbb{R}\) and any \(\varepsilon > 0\) we have \(B(g,\varepsilon) \cap \emptyset = \emptyset\), so no \(g\) satisfies the defining condition 6.7. Thus

\begin{equation*} \overline{U \cap W} = \emptyset . \end{equation*}

Second, \(1 \in \overline{U} \cap \overline{W}\). To see that \(1 \in \overline{U}\), let \(\varepsilon > 0\) and put

\begin{equation*} t = \max\left\{\tfrac{1}{2},\, 1 - \tfrac{\varepsilon}{2}\right\}. \end{equation*}

Then \(t \geq 1/2 > 0\) and \(t < 1\) (since both \(1/2 < 1\) and \(1 - \varepsilon/2 < 1\)), so \(t \in (0,1) = U\); moreover \(1 - t \leq \varepsilon/2 < \varepsilon\), so \(t \in B(1,\varepsilon)\). Hence \(B(1,\varepsilon) \cap U \neq \emptyset\) for every \(\varepsilon > 0\), so \(1 \in \overline{U}\).

Symmetrically, to see that \(1 \in \overline{W}\), let \(\varepsilon > 0\) and put \(s = \min\{3/2,\, 1 + \varepsilon/2\}\). Then \(s > 1\) and \(s \leq 3/2 < 2\), so \(s \in (1,2) = W\), while \(s - 1 \leq \varepsilon/2 < \varepsilon\), so \(s \in B(1,\varepsilon)\). Hence \(1 \in \overline{W}\).

Therefore \(1 \in \overline{U} \cap \overline{W}\) but \(1 \notin \overline{U \cap W} = \emptyset\), so

\begin{equation*} \overline{U} \cap \overline{W} \neq \overline{U \cap W}, \end{equation*}

and the proposed statement fails.

An even more dramatic counterexample in \(V = \mathbb{R}\): take \(U = \mathbb{Q}\) and \(W = \mathbb{R} \setminus \mathbb{Q}\). Every open interval contains both a rational and an irrational number, so \(\overline{U} = \overline{W} = \mathbb{R}\) and hence \(\overline{U} \cap \overline{W} = \mathbb{R}\), while \(U \cap W = \emptyset\) gives \(\overline{U \cap W} = \emptyset\).

One inclusion does always hold. The closure operation is order preserving: if \(E_1 \subseteq E_2 \subseteq V\) and \(g \in \overline{E_1}\), then for every \(\varepsilon > 0\) we have \(\emptyset \neq B(g,\varepsilon) \cap E_1 \subseteq B(g,\varepsilon) \cap E_2\), so \(g \in \overline{E_2}\). Since \(U \cap W \subseteq U\) and \(U \cap W \subseteq W\), this gives \(\overline{U \cap W} \subseteq \overline{U}\) and \(\overline{U \cap W} \subseteq \overline{W}\), hence

\begin{equation*} \overline{U \cap W} \subseteq \overline{U} \cap \overline{W}. \end{equation*}

The counterexamples above show that this inclusion can be strict.

Problem (6A.12)

Suppose \((U, d_U)\), \((V, d_V)\), and \((W, d_W)\) are metric spaces. Suppose also that \(T : U \to V\) and \(S : V \to W\) are continuous functions.

(a) Using the definition of continuity, show that \(S \circ T : U \to W\) is continuous.

(b) Using the equivalence of 6.11(a) and 6.11(b), show that \(S \circ T : U \to W\) is continuous.

(c) Using the equivalence of 6.11(a) and 6.11(c), show that \(S \circ T : U \to W\) is continuous.

Solution

(a) We use definition 6.10 directly. Fix \(f \in U\) and let \(\varepsilon > 0\).

Because \(S\) is continuous, \(S\) is continuous at the element \(T(f) \in V\), so there exists \(\gamma > 0\) such that

\begin{equation*} d_W\bigl(S(T(f)), S(v)\bigr) < \varepsilon \quad \text{for all } v \in V \text{ with } d_V\bigl(T(f), v\bigr) < \gamma . \end{equation*}

Because \(T\) is continuous, \(T\) is continuous at \(f\), so with this \(\gamma\) in the role of the tolerance there exists \(\delta > 0\) such that

\begin{equation*} d_V\bigl(T(f), T(g)\bigr) < \gamma \quad \text{for all } g \in U \text{ with } d_U(f,g) < \delta . \end{equation*}

Now suppose \(g \in U\) satisfies \(d_U(f,g) < \delta\). Then \(d_V(T(f), T(g)) < \gamma\), and applying the first display with \(v = T(g)\) gives

\begin{equation*} d_W\bigl((S \circ T)(f), (S \circ T)(g)\bigr) = d_W\bigl(S(T(f)), S(T(g))\bigr) < \varepsilon . \end{equation*}

Thus \(S \circ T\) is continuous at \(f\). Since \(f \in U\) was arbitrary, \(S \circ T\) is continuous.

(b) We verify condition 6.11(b) for \(S \circ T\). Suppose \(f_1, f_2, \ldots\) is a sequence in \(U\) with \(\lim_{k \to \infty} f_k = f\) for some \(f \in U\).

Because \(T\) is continuous, the implication 6.11(a) \(\Rightarrow\) 6.11(b) applied to \(T\) gives

\begin{equation*} \lim_{k \to \infty} T(f_k) = T(f) \quad \text{in } V . \end{equation*}

Because \(S\) is continuous, the implication 6.11(a) \(\Rightarrow\) 6.11(b) applied to \(S\), with the convergent sequence \(T(f_1), T(f_2), \ldots\) in \(V\), gives

\begin{equation*} \lim_{k \to \infty} S\bigl(T(f_k)\bigr) = S\bigl(T(f)\bigr) \quad \text{in } W , \end{equation*}

that is, \(\lim_{k \to \infty} (S \circ T)(f_k) = (S \circ T)(f)\).

Thus \(S \circ T\) satisfies 6.11(b), and therefore the implication 6.11(b) \(\Rightarrow\) 6.11(a) shows that \(S \circ T\) is continuous.

(c) We verify condition 6.11(c) for \(S \circ T\). First note that for every set \(G \subseteq W\),

\begin{equation*} (S \circ T)^{-1}(G) = T^{-1}\bigl(S^{-1}(G)\bigr), \end{equation*}

because for \(f \in U\) we have the chain of equivalences: \(f \in (S \circ T)^{-1}(G)\) if and only if \(S(T(f)) \in G\), if and only if \(T(f) \in S^{-1}(G)\), if and only if \(f \in T^{-1}(S^{-1}(G))\).

Now suppose \(G\) is an open subset of \(W\). Because \(S\) is continuous, the implication 6.11(a) \(\Rightarrow\) 6.11(c) applied to \(S\) shows that \(S^{-1}(G)\) is an open subset of \(V\). Because \(T\) is continuous, the same implication applied to \(T\), with the open set \(S^{-1}(G) \subseteq V\), shows that \(T^{-1}(S^{-1}(G))\) is an open subset of \(U\). By the displayed identity, \((S \circ T)^{-1}(G)\) is an open subset of \(U\).

Thus \(S \circ T\) satisfies 6.11(c), and therefore the implication 6.11(c) \(\Rightarrow\) 6.11(a) shows that \(S \circ T\) is continuous.

Problem (6A.13)

Prove the parts of 6.11 that were not proved in the text.

[6.11 states: Suppose \(V\) and \(W\) are metric spaces and \(T : V \to W\) is a function. Then the following are equivalent.

(a) \(T\) is continuous.

(b) \(\lim_{k \to \infty} f_k = f\) in \(V\) implies \(\lim_{k \to \infty} T(f_k) = T(f)\) in \(W\).

(c) \(T^{-1}(G)\) is an open subset of \(V\) for every open set \(G \subseteq W\).

(d) \(T^{-1}(F)\) is a closed subset of \(V\) for every closed set \(F \subseteq W\).

The text proves that (b) implies (d), and that (c) and (d) are equivalent.]

Solution

Write \(d_V\) and \(d_W\) for the metrics on \(V\) and \(W\). The text proves the implication (b) \(\Rightarrow\) (d) and the equivalence (c) \(\Leftrightarrow\) (d). To complete the proof that all four conditions are equivalent, it suffices to prove (a) \(\Rightarrow\) (b) and (c) \(\Rightarrow\) (a). Indeed, these together with the text’s results give the cycle

\begin{equation*} \text{(a)} \Rightarrow \text{(b)} \Rightarrow \text{(d)} \Rightarrow \text{( c)} \Rightarrow \text{(a)}, \end{equation*}

from which every one of the four conditions implies every other.

Proof that (a) implies (b). Suppose \(T\) is continuous. Suppose \(f_1, f_2, \ldots\) is a sequence in \(V\) with \(\lim_{k \to \infty} f_k = f\) for some \(f \in V\). Let \(\varepsilon > 0\). Because \(T\) is continuous at \(f\), definition 6.10 provides \(\delta > 0\) such that

\begin{equation*} d_W\bigl(T(f), T(g)\bigr) < \varepsilon \quad \text{for all } g \in V \text{ with } d_V(f,g) < \delta . \end{equation*}

Because \(\lim_{k \to \infty} f_k = f\) means \(\lim_{k \to \infty} d_V(f_k, f) = 0\) (definition 6.8), there exists \(n \in \mathbb{Z}^+\) such that \(d_V(f_k, f) < \delta\) for all integers \(k \geq n\). For such \(k\) the display above (with \(g = f_k\)) gives \(d_W(T(f), T(f_k)) < \varepsilon\).

Thus for every \(\varepsilon > 0\) there exists \(n \in \mathbb{Z}^+\) with \(d_W(T(f_k), T(f)) < \varepsilon\) for all \(k \geq n\); that is, \(\lim_{k \to \infty} d_W(T(f_k), T(f)) = 0\), which means \(\lim_{k \to \infty} T(f_k) = T(f)\). Hence (b) holds.

Proof that (c) implies (a). Suppose \(T^{-1}(G)\) is an open subset of \(V\) for every open set \(G \subseteq W\). Let \(f \in V\) and let \(\varepsilon > 0\).

The open ball \(B\bigl(T(f), \varepsilon\bigr)\) is an open subset of \(W\) by 6.5. Hence

\begin{equation*} G := T^{-1}\Bigl(B\bigl(T(f), \varepsilon\bigr)\Bigr) \end{equation*}

is an open subset of \(V\). Since \(d_W(T(f), T(f)) = 0 < \varepsilon\), we have \(T(f) \in B(T(f),\varepsilon)\) and therefore \(f \in G\). Because \(G\) is open, definition 6.4 provides \(\delta > 0\) with \(B(f,\delta) \subseteq G\).

Now suppose \(g \in V\) satisfies \(d_V(f,g) < \delta\). Then \(g \in B(f,\delta) \subseteq G\), so \(T(g) \in B(T(f),\varepsilon)\), which means

\begin{equation*} d_W\bigl(T(f), T(g)\bigr) < \varepsilon . \end{equation*}

Thus \(T\) is continuous at \(f\). Since \(f \in V\) was arbitrary, \(T\) is continuous, so (a) holds.

For completeness, here is a direct proof of the implication (a) \(\Rightarrow\) (c), which shows the equivalence without routing through (b) and (d). Suppose \(T\) is continuous and \(G \subseteq W\) is open. Let \(f \in T^{-1}(G)\). Then \(T(f) \in G\), so there exists \(\varepsilon > 0\) with \(B(T(f),\varepsilon) \subseteq G\). By continuity of \(T\) at \(f\), there exists \(\delta > 0\) such that \(d_W(T(f), T(g)) < \varepsilon\) whenever \(d_V(f,g) < \delta\). Hence \(g \in B(f,\delta)\) implies \(T(g) \in B(T(f),\varepsilon) \subseteq G\), that is, \(B(f,\delta) \subseteq T^{-1}(G)\). Therefore \(T^{-1}(G)\) is open.

Problem (6A.14)

Suppose a Cauchy sequence in a metric space has a convergent subsequence. Prove that the Cauchy sequence converges.

Solution

Suppose \((V,d)\) is a metric space and \(f_1, f_2, \ldots\) is a Cauchy sequence in \(V\). Suppose \(k_1 < k_2 < \cdots\) are positive integers and \(f \in V\) is such that

\begin{equation*} \lim_{m \to \infty} f_{k_m} = f . \end{equation*}

We show that \(\lim_{k \to \infty} f_k = f\).

First we record that \(k_m \geq m\) for every \(m \in \mathbb{Z}^+\). This holds by induction: \(k_1 \geq 1\) because \(k_1\) is a positive integer, and if \(k_m \geq m\) then \(k_{m+1} > k_m \geq m\), so \(k_{m+1} \geq m+1\) because \(k_{m+1}\) is an integer.

Now let \(\varepsilon > 0\). Because \(f_1, f_2, \ldots\) is a Cauchy sequence (definition 6.12), there exists \(n \in \mathbb{Z}^+\) such that

\begin{equation*} d(f_j, f_k) < \frac{\varepsilon}{2} \quad \text{for all integers } j \geq n \text{ and } k \geq n . \end{equation*}

Because \(\lim_{m \to \infty} f_{k_m} = f\), which by 6.8 means \(\lim_{m \to \infty} d(f_{k_m}, f) = 0\), there exists \(M \in \mathbb{Z}^+\) such that \(d(f_{k_m}, f) < \varepsilon/2\) for all \(m \geq M\). Choose an integer \(m \geq \max\{M, n\}\) and set \(j = k_m\). Then

\begin{equation*} d(f_j, f) = d(f_{k_m}, f) < \frac{\varepsilon}{2}, \end{equation*}

and also \(j = k_m \geq m \geq n\) by the observation of the previous paragraph.

Now suppose \(k \in \mathbb{Z}^+\) with \(k \geq n\). Since \(j \geq n\) as well, the Cauchy estimate gives \(d(f_k, f_j) < \varepsilon/2\), and the triangle inequality then gives

\begin{equation*} d(f_k, f) \leq d(f_k, f_j) + d(f_j, f) < \frac{\varepsilon}{2} + \frac{\varepsilon}{2} = \varepsilon . \end{equation*}

Thus for every \(\varepsilon > 0\) there exists \(n \in \mathbb{Z}^+\) such that \(d(f_k, f) < \varepsilon\) for all integers \(k \geq n\). In other words \(\lim_{k \to \infty} d(f_k, f) = 0\), which by 6.8 means \(\lim_{k \to \infty} f_k = f\). Hence the Cauchy sequence converges (to the same limit as its convergent subsequence).

Problem (6A.15)

Verify that all five of the metric spaces in Example 6.2 are complete metric spaces.

Solution

Recall the five metric spaces of 6.2: (i) an arbitrary nonempty set \(V\) with \(d(f,g) = 1\) if \(f \neq g\) and \(d(f,f) = 0\); (ii) \(\mathbb{R}\) with \(d(x,y) = |x-y|\); (iii) \(\mathbb{R}^n\) with \(d\big((x_1,\dots,x_n),(y_1,\dots,y_n)\big) = \max\{|x_1-y_1|,\dots,|x_n-y_n|\}\); (iv) \(C([0,1])\), the continuous real-valued functions on \([0,1]\), with \(d(f,g) = \sup\{|f(t)-g(t)| : t \in [0,1]\}\); (v) \(\ell^1\), the set of real sequences \((a_1,a_2,\dots)\) with \(\sum_{k=1}^{\infty}|a_k| < \infty\), with \(d\big((a_1,a_2,\dots),(b_1,b_2,\dots)\big) = \sum_{k=1}^{\infty}|a_k-b_k|\).

(i) The discrete metric. Suppose \(f_1, f_2, \dots\) is a Cauchy sequence in \(V\). Applying the definition 6.12 with \(\varepsilon = 1\), there exists \(n \in \mathbb{Z}^+\) such that \(d(f_j,f_k) < 1\) for all \(j,k \geq n\). Because \(d\) takes only the values \(0\) and \(1\), this forces \(d(f_j,f_k) = 0\), and hence \(f_j = f_k\), for all \(j,k \geq n\). Thus \(f_k = f_n\) for every \(k \geq n\), so \(d(f_k, f_n) = 0 < \varepsilon\) for every \(k \geq n\) and every \(\varepsilon > 0\). Hence \(\lim_{k \to \infty} f_k = f_n \in V\), and \(V\) is complete.

(ii) \(\mathbb{R}\) with the absolute-value metric. Suppose \(x_1, x_2, \dots\) is a Cauchy sequence of real numbers. First, the sequence is bounded: taking \(\varepsilon = 1\) in 6.12 gives \(n \in \mathbb{Z}^+\) with \(|x_j - x_k| < 1\) for all \(j,k \geq n\); hence \(|x_k| \leq |x_n| + 1\) for all \(k \geq n\), and therefore

\begin{equation*} |x_k| \leq \max\{|x_1|, \dots, |x_{n-1}|, |x_n| + 1\} \end{equation*}

for all \(k \in \mathbb{Z}^+\). Let \(M\) denote the right side above. For each \(m \in \mathbb{Z}^+\) the set \(\{x_k : k \geq m\}\) is nonempty and bounded above by \(M\), so by the completeness (least upper bound) property of \(\mathbb{R}\) the number

\begin{equation*} s_m = \sup\{x_k : k \geq m\} \end{equation*}

exists; furthermore \(-M \leq s_m\) and \(s_1 \geq s_2 \geq \cdots\). Hence \(\{s_m : m \in \mathbb{Z}^+\}\) is a nonempty set of real numbers bounded below by \(-M\), so \(L = \inf\{s_m : m \in \mathbb{Z}^+\}\) exists in \(\mathbb{R}\).

We claim \(\lim_{k \to \infty} x_k = L\). Suppose \(\varepsilon > 0\). There exists \(n \in \mathbb{Z}^+\) such that \(|x_j - x_k| < \varepsilon/2\) for all \(j,k \geq n\). For every \(k \geq n\) we have \(x_k \leq x_n + \varepsilon/2\); taking the supremum over \(k \geq n\) gives \(s_n \leq x_n + \varepsilon/2\), and thus \(L \leq s_n \leq x_n + \varepsilon/2\). On the other hand, for every \(m \geq n\) we have \(s_m \geq x_m > x_n - \varepsilon/2\); because \(s_1 \geq s_2 \geq \cdots\) we have \(L = \inf\{s_m : m \geq n\} \geq x_n - \varepsilon/2\). Combining the two estimates gives \(|x_n - L| \leq \varepsilon/2\). Hence for every \(j \geq n\),

\begin{equation*} |x_j - L| \leq |x_j - x_n| + |x_n - L| < \tfrac{\varepsilon}{2} + \tfrac{\varepsilon}{2} = \varepsilon . \end{equation*}

Thus \(x_k \to L \in \mathbb{R}\), and \(\mathbb{R}\) is complete.

(iii) \(\mathbb{R}^n\) with the maximum metric. Suppose \(x^{(1)}, x^{(2)}, \dots\) is a Cauchy sequence in \(\mathbb{R}^n\), where \(x^{(k)} = (x^{(k)}_1, \dots, x^{(k)}_n)\). For each \(i \in \{1,\dots,n\}\) and all \(j,k\) we have

\begin{equation*} |x^{(j)}_i - x^{(k)}_i| \leq \max_{1 \leq p \leq n} |x^{(j)}_p - x^{(k)}_p| = d\big(x^{(j)}, x^{(k)}\big), \end{equation*}

so the real sequence \(x^{(1)}_i, x^{(2)}_i, \dots\) is Cauchy. By (ii) it converges to some \(y_i \in \mathbb{R}\). Let \(y = (y_1, \dots, y_n)\). Suppose \(\varepsilon > 0\). For each \(i \in \{1,\dots,n\}\) there exists \(m_i \in \mathbb{Z}^+\) such that \(|x^{(k)}_i - y_i| < \varepsilon\) for all \(k \geq m_i\). Because there are only finitely many coordinates, \(N = \max\{m_1, \dots, m_n\}\) is a positive integer, and for \(k \geq N\) we have

\begin{equation*} d\big(x^{(k)}, y\big) = \max_{1 \leq i \leq n} |x^{(k)}_i - y_i| < \varepsilon . \end{equation*}

Thus \(x^{(k)} \to y \in \mathbb{R}^n\), and \(\mathbb{R}^n\) is complete. (Finiteness of the number of coordinates is essential here; it is what allows the single \(N\).)

(iv) \(C([0,1])\) with the supremum metric. Suppose \(f_1, f_2, \dots\) is a Cauchy sequence in \(C([0,1])\). For each \(t \in [0,1]\) and all \(j,k\) we have \(|f_j(t) - f_k(t)| \leq d(f_j,f_k)\), so \(f_1(t), f_2(t), \dots\) is a Cauchy sequence of real numbers; by (ii) it converges. Define \(f : [0,1] \to \mathbb{R}\) by

\begin{equation*} f(t) = \lim_{k \to \infty} f_k(t). \end{equation*}

We first show that \(\sup\{|f(t) - f_k(t)| : t \in [0,1]\} \to 0\). Suppose \(\varepsilon > 0\). There exists \(n \in \mathbb{Z}^+\) such that \(d(f_j,f_k) < \varepsilon/3\) for all \(j,k \geq n\). Fix \(k \geq n\) and \(t \in [0,1]\). Then \(|f_j(t) - f_k(t)| < \varepsilon/3\) for all \(j \geq n\); letting \(j \to \infty\) and using the continuity of the absolute value on \(\mathbb{R}\) gives

\begin{equation*} |f(t) - f_k(t)| \leq \tfrac{\varepsilon}{3} \end{equation*}

for all \(t \in [0,1]\) and all \(k \geq n\). Hence \(\sup\{|f(t) - f_k(t)| : t \in [0,1]\} \leq \varepsilon/3\) for \(k \geq n\).

Next we show \(f\) is continuous. Suppose \(t \in [0,1]\) and \(\varepsilon > 0\); take \(n\) as above and set \(k = n\). Because \(f_n\) is continuous at \(t\), there exists \(\delta > 0\) such that \(|f_n(t) - f_n(s)| < \varepsilon/3\) for all \(s \in [0,1]\) with \(|s - t| < \delta\). For such \(s\),

\begin{equation*} |f(t) - f(s)| \leq |f(t) - f_n(t)| + |f_n(t) - f_n(s)| + |f_n(s) - f(s)| < \tfrac{\varepsilon}{3} + \tfrac{\varepsilon}{3} + \tfrac{\varepsilon}{3} = \varepsilon . \end{equation*}

Thus \(f\) is continuous at each \(t \in [0,1]\), so \(f \in C([0,1])\). Now \(d(f_k, f) = \sup\{|f(t)-f_k(t)| : t \in [0,1]\} \leq \varepsilon/3 < \varepsilon\) for all \(k \geq n\), so \(\lim_{k \to \infty} f_k = f\) in \(C([0,1])\). Hence \(C([0,1])\) is complete.

(v) \(\ell^1\). Suppose \(a^{(1)}, a^{(2)}, \dots\) is a Cauchy sequence in \(\ell^1\), where \(a^{(m)} = (a^{(m)}_1, a^{(m)}_2, \dots)\). For each \(k \in \mathbb{Z}^+\) and all \(j,m\),

\begin{equation*} |a^{(j)}_k - a^{(m)}_k| \leq \sum_{p=1}^{\infty} |a^{(j)}_p - a^{(m)}_p| = d\big(a^{(j)}, a^{(m)}\big), \end{equation*}

so the real sequence \(a^{(1)}_k, a^{(2)}_k, \dots\) is Cauchy; by (ii) it converges to some \(b_k \in \mathbb{R}\). Let \(b = (b_1, b_2, \dots)\).

Suppose \(\varepsilon > 0\). There exists \(n \in \mathbb{Z}^+\) such that \(d(a^{(j)}, a^{(m)}) < \varepsilon/2\) for all \(j,m \geq n\). Fix \(m \geq n\) and \(N \in \mathbb{Z}^+\). For every \(j \geq n\),

\begin{equation*} \sum_{k=1}^{N} |a^{(j)}_k - a^{(m)}_k| \leq d\big(a^{(j)}, a^{(m)}\big) \leq \tfrac{\varepsilon}{2}. \end{equation*}

The left side is a finite sum, so letting \(j \to \infty\) (each of the \(N\) terms converges) gives

\begin{equation*} \sum_{k=1}^{N} |b_k - a^{(m)}_k| \leq \tfrac{\varepsilon}{2}. \end{equation*}

This holds for every \(N \in \mathbb{Z}^+\); the partial sums of the nonnegative series \(\sum_k |b_k - a^{(m)}_k|\) are therefore bounded above by \(\varepsilon/2\), and hence

\begin{equation*} \sum_{k=1}^{\infty} |b_k - a^{(m)}_k| \leq \tfrac{\varepsilon}{2} < \varepsilon \end{equation*}

for all \(m \geq n\).

Taking \(m = n\), the sequence \(b - a^{(n)}\) lies in \(\ell^1\). Because \(a^{(n)} \in \ell^1\) and

\begin{equation*} \sum_{k=1}^{\infty} |b_k| \leq \sum_{k=1}^{\infty} |b_k - a^{(n)}_k| + \sum_{k=1}^{\infty} |a^{(n)}_k| < \infty, \end{equation*}

we conclude \(b \in \ell^1\). The displayed estimate now reads \(d(a^{(m)}, b) < \varepsilon\) for all \(m \geq n\), so \(\lim_{m \to \infty} a^{(m)} = b\) in \(\ell^1\). Hence \(\ell^1\) is complete.

Problem (6A.16)

Suppose \((U, d)\) is a metric space. Let \(W\) denote the set of all Cauchy sequences of elements of \(U\).

(a) For \((f_1, f_2, \dots)\) and \((g_1, g_2, \dots)\) in \(W\), define \((f_1, f_2, \dots) \equiv (g_1, g_2, \dots)\) to mean that

\begin{equation*} \lim_{k \to \infty} d(f_k, g_k) = 0 . \end{equation*}

Show that \(\equiv\) is an equivalence relation on \(W\).

(b) Let \(V\) denote the set of equivalence classes of elements of \(W\) under the equivalence relation above. For \((f_1, f_2, \dots) \in W\), let \(\widehat{(f_1, f_2, \dots)}\) denote the equivalence class of \((f_1, f_2, \dots)\). Define \(d_V : V \times V \to [0, \infty)\) by

\begin{equation*} d_V\big(\widehat{(f_1, f_2, \dots)}, \widehat{(g_1, g_2, \dots)}\big) = \lim_{k \to \infty} d(f_k, g_k) . \end{equation*}

Show that this definition of \(d_V\) makes sense and that \(d_V\) is a metric on \(V\).

(c) Show that \((V, d_V)\) is a complete metric space.

(d) Show that the map from \(U\) to \(V\) that takes \(f \in U\) to \(\widehat{(f, f, f, \dots)}\) preserves distances, meaning that

\begin{equation*} d(f, g) = d_V\big(\widehat{(f, f, f, \dots)}, \widehat{(g, g, g, \dots)}\big) \end{equation*}

for all \(f, g \in U\).

(e) Explain why (d) shows that every metric space is a subset of some complete metric space.

Solution

Throughout we use freely the fact that \(\mathbb{R}\) with the metric \(|x-y|\) is complete (Example 6.15, verified in Exercise 15 of this section).

(a) Reflexivity. If \((f_1,f_2,\dots) \in W\), then \(d(f_k,f_k) = 0\) for every \(k\) by the first bullet point of 6.1, so \(\lim_{k \to \infty} d(f_k,f_k) = 0\) and \((f_1,f_2,\dots) \equiv (f_1,f_2,\dots)\).

Symmetry. Suppose \((f_1,f_2,\dots) \equiv (g_1,g_2,\dots)\). Because \(d(g_k,f_k) = d(f_k,g_k)\) for every \(k\) (third bullet point of 6.1), we get \(\lim_{k \to \infty} d(g_k,f_k) = 0\), so \((g_1,g_2,\dots) \equiv (f_1,f_2,\dots)\).

Transitivity. Suppose \((f_1,f_2,\dots) \equiv (g_1,g_2,\dots)\) and \((g_1,g_2,\dots) \equiv (h_1,h_2,\dots)\). The triangle inequality (fourth bullet point of 6.1) gives

\begin{equation*} 0 \leq d(f_k,h_k) \leq d(f_k,g_k) + d(g_k,h_k) \end{equation*}

for every \(k \in \mathbb{Z}^+\). Both terms on the right tend to \(0\), so \(\lim_{k \to \infty} d(f_k,h_k) = 0\) and \((f_1,f_2,\dots) \equiv (h_1,h_2,\dots)\).

Thus \(\equiv\) is an equivalence relation on \(W\).

(b) Two things must be checked: that the limit defining \(d_V\) exists, and that its value does not depend on the representatives chosen from the two equivalence classes.

The key inequality. Suppose \(f, f^{\prime}, g, g^{\prime} \in U\). Two applications of the triangle inequality give

\begin{equation*} d(f,g) \leq d(f,f^{\prime}) + d(f^{\prime},g^{\prime}) + d(g^{\prime},g), \end{equation*}

so \(d(f,g) - d(f^{\prime},g^{\prime}) \leq d(f,f^{\prime}) + d(g,g^{\prime})\). Interchanging the roles of the primed and unprimed elements gives the reverse inequality, and hence

\begin{equation*} \big| d(f,g) - d(f^{\prime},g^{\prime}) \big| \leq d(f,f^{\prime}) + d(g,g^{\prime}) . \end{equation*}

Existence of the limit. Suppose \((f_1,f_2,\dots)\) and \((g_1,g_2,\dots)\) are in \(W\). Applying the key inequality with \(f = f_j\), \(f^{\prime} = f_k\), \(g = g_j\), \(g^{\prime} = g_k\) gives

\begin{equation*} \big| d(f_j,g_j) - d(f_k,g_k) \big| \leq d(f_j,f_k) + d(g_j,g_k) . \end{equation*}

Suppose \(\varepsilon > 0\). Because both sequences are Cauchy, there exists \(n \in \mathbb{Z}^+\) such that \(d(f_j,f_k) < \varepsilon/2\) and \(d(g_j,g_k) < \varepsilon/2\) for all \(j,k \geq n\); the displayed inequality then gives \(|d(f_j,g_j) - d(f_k,g_k)| < \varepsilon\) for all \(j,k \geq n\). Thus \(d(f_1,g_1), d(f_2,g_2), \dots\) is a Cauchy sequence of real numbers, so it converges to some real number. Each term is nonnegative, so the limit lies in \([0,\infty)\), as required by the stated codomain of \(d_V\).

Independence of representatives. Suppose \((f_1,f_2,\dots) \equiv (f^{\prime}_1,f^{\prime}_2,\dots)\) and \((g_1,g_2,\dots) \equiv (g^{\prime}_1,g^{\prime}_2,\dots)\), all four sequences being in \(W\). Applying the key inequality with \(f = f_k\), \(f^{\prime} = f^{\prime}_k\), \(g = g_k\), \(g^{\prime} = g^{\prime}_k\) gives

\begin{equation*} \big| d(f_k,g_k) - d(f^{\prime}_k,g^{\prime}_k) \big| \leq d(f_k,f^{\prime}_k) + d(g_k,g^{\prime}_k) \to 0 \end{equation*}

as \(k \to \infty\). Because both limits exist by the previous paragraph, we conclude

\begin{equation*} \lim_{k \to \infty} d(f_k,g_k) = \lim_{k \to \infty} d(f^{\prime}_k,g^{\prime}_k) . \end{equation*}

Hence the value assigned by \(d_V\) depends only on the two equivalence classes, and \(d_V : V \times V \to [0,\infty)\) is a well-defined function.

\(d_V\) is a metric. Write \(F = \widehat{(f_1,f_2,\dots)}\), \(G = \widehat{(g_1,g_2,\dots)}\), \(H = \widehat{(h_1,h_2,\dots)}\) for arbitrary elements of \(V\).

If \(F \in V\), then \(d_V(F,F) = \lim_{k \to \infty} d(f_k,f_k) = 0\).

If \(d_V(F,G) = 0\), then \(\lim_{k \to \infty} d(f_k,g_k) = 0\), which says exactly that \((f_1,f_2,\dots) \equiv (g_1,g_2,\dots)\); hence these two sequences lie in the same equivalence class, which means \(F = G\).

Symmetry: \(d_V(F,G) = \lim_{k \to \infty} d(f_k,g_k) = \lim_{k \to \infty} d(g_k,f_k) = d_V(G,F)\).

Triangle inequality: for every \(k\) we have \(d(f_k,h_k) \leq d(f_k,g_k) + d(g_k,h_k)\). All three limits exist by the first part of (b), so taking limits as \(k \to \infty\) preserves the inequality:

\begin{equation*} d_V(F,H) \leq d_V(F,G) + d_V(G,H) . \end{equation*}

Thus \(d_V\) satisfies all four bullet points of 6.1, and \((V, d_V)\) is a metric space. (Note that \(V\) is nonempty because \(U\) is nonempty: any constant sequence is Cauchy.)

(c) We begin by recording the constant-sequence embedding, which is also the content of (d) and which will be used several times below. Suppose \(f \in U\). The constant sequence \((f,f,f,\dots)\) satisfies \(d(f,f) = 0 < \varepsilon\) for every \(\varepsilon > 0\), so it is a Cauchy sequence and thus lies in \(W\). Define \(\iota : U \to V\) by

\begin{equation*} \iota(f) = \widehat{(f,f,f,\dots)} . \end{equation*}

For all \(f,g \in U\) the sequence whose \(k\)-th term is \(d(f,g)\) is constant, so

\begin{equation*} d_V\big(\iota(f), \iota(g)\big) = \lim_{k \to \infty} d(f,g) = d(f,g) . \end{equation*}

Now suppose \(F^{(1)}, F^{(2)}, \dots\) is a Cauchy sequence in \((V, d_V)\). For each \(m \in \mathbb{Z}^+\) choose a representative, so that \(F^{(m)} = \widehat{(f^{(m)}_1, f^{(m)}_2, \dots)}\) with \((f^{(m)}_1, f^{(m)}_2, \dots) \in W\). Because this sequence is Cauchy in \(U\), there exists \(k_m \in \mathbb{Z}^+\) such that

\begin{equation*} d\big(f^{(m)}_j, f^{(m)}_k\big) < \tfrac{1}{m} \quad \text{for all } j, k \geq k_m . \end{equation*}

Set \(h_m = f^{(m)}_{k_m}\). Then for every \(k \geq k_m\) we have \(d(f^{(m)}_k, h_m) < 1/m\), and therefore

\begin{equation*} d_V\big(F^{(m)}, \iota(h_m)\big) = \lim_{k \to \infty} d\big(f^{(m)}_k, h_m\big) \leq \tfrac{1}{m} . \end{equation*}

The sequence \(h_1, h_2, \dots\) is Cauchy in \(U\). Using the displayed identity \(d(h_j,h_m) = d_V(\iota(h_j), \iota(h_m))\) and then the triangle inequality in \(V\),

\begin{equation*} \begin{aligned} d(h_j, h_m) &= d_V\big(\iota(h_j), \iota(h_m)\big) \\ &\leq d_V\big(\iota(h_j), F^{(j)}\big) + d_V\big(F^{(j)}, F^{(m)}\big) + d_V\big(F^{(m)}, \iota(h_m)\big) \\ &\leq \tfrac{1}{j} + d_V\big(F^{(j)}, F^{(m)}\big) + \tfrac{1}{m}, \end{aligned} \end{equation*}

where the last line uses the displayed bound \(d_V(F^{(m)}, \iota(h_m)) \leq 1/m\). Suppose \(\varepsilon > 0\). Because \(F^{(1)}, F^{(2)}, \dots\) is Cauchy in \(V\), there exists \(n_0 \in \mathbb{Z}^+\) such that \(d_V(F^{(j)}, F^{(m)}) < \varepsilon/3\) for all \(j,m \geq n_0\). Let \(n\) be an integer with \(n \geq n_0\) and \(n > 3/\varepsilon\). Then for all \(j,m \geq n\),

\begin{equation*} d(h_j,h_m) < \tfrac{\varepsilon}{3} + \tfrac{\varepsilon}{3} + \tfrac{\varepsilon}{3} = \varepsilon . \end{equation*}

Thus \((h_1,h_2,\dots) \in W\); let \(H = \widehat{(h_1,h_2,\dots)} \in V\).

\(F^{(m)} \to H\). First note that for each \(m\),

\begin{equation*} d_V\big(\iota(h_m), H\big) = \lim_{k \to \infty} d(h_m, h_k). \end{equation*}

Suppose \(\varepsilon > 0\). Because \(h_1, h_2, \dots\) is Cauchy, there exists \(n_1 \in \mathbb{Z}^+\) such that \(d(h_m, h_k) < \varepsilon/3\) for all \(m,k \geq n_1\); letting \(k \to \infty\) gives \(d_V(\iota(h_m), H) \leq \varepsilon/3\) for all \(m \geq n_1\). Now let \(n\) be an integer with \(n \geq n_1\) and \(n > 3/\varepsilon\). For all \(m \geq n\), using \(d_V(F^{(m)}, \iota(h_m)) \leq 1/m\) and the triangle inequality in \(V\),

\begin{equation*} d_V\big(F^{(m)}, H\big) \leq d_V\big(F^{(m)}, \iota(h_m)\big) + d_V\big(\iota(h_m), H\big) \leq \tfrac{1}{m} + \tfrac{\varepsilon}{3} < \tfrac{\varepsilon}{3} + \tfrac{\varepsilon}{3} < \varepsilon . \end{equation*}

Hence \(\lim_{m \to \infty} F^{(m)} = H\) in \((V, d_V)\). Every Cauchy sequence in \(V\) converges to an element of \(V\), so \((V,d_V)\) is a complete metric space by 6.14.

(d) This is the identity recorded at the start of (c), which we now restate. For \(f \in U\) the constant sequence \((f,f,f,\dots)\) is Cauchy and hence belongs to \(W\), so the map \(\iota : U \to V\) given by \(\iota(f) = \widehat{(f,f,f,\dots)}\) is defined. If \(f, g \in U\), then \(d(f_k, g_k) = d(f,g)\) for every \(k\) when \(f_k = f\) and \(g_k = g\) for all \(k\); this constant sequence of real numbers has limit \(d(f,g)\), and therefore

\begin{equation*} d_V\big(\widehat{(f,f,f,\dots)}, \widehat{(g,g,g,\dots)}\big) = \lim_{k \to \infty} d(f,g) = d(f,g) . \end{equation*}

Thus \(\iota\) preserves distances. (This computation uses only the definition of \(d_V\), which part (b) showed to be independent of the representatives chosen, so no circularity arises from having invoked it inside (c).)

(e) By (d), the map \(\iota : U \to V\) preserves distances. In particular \(\iota\) is injective: if \(\iota(f) = \iota(g)\), then \(d(f,g) = d_V(\iota(f), \iota(g)) = d_V(\iota(f), \iota(f)) = 0\), and hence \(f = g\) by the second bullet point of 6.1. Thus \(\iota\) is a bijection from \(U\) onto the subset \(\iota(U)\) of \(V\), and it carries the metric \(d\) to the metric that \(\iota(U)\) inherits from \((V,d_V)\): for \(f,g \in U\),

\begin{equation*} d(f,g) = d_V\big(\iota(f), \iota(g)\big). \end{equation*}

So \((U,d)\) and \((\iota(U), d_V|_{\iota(U) \times \iota(U)})\) are indistinguishable as metric spaces; every metric statement about one is the same statement about the other, with \(f\) replaced by \(\iota(f)\). Identifying each \(f \in U\) with \(\iota(f)\), we may regard \(U\) itself as a subset of \(V\), carrying its original metric. Since \((V,d_V)\) is complete by (c), every metric space is (after this identification) a subset of a complete metric space.

Stated without the identification: every metric space is isometric to a subset of a complete metric space. One can say more, though it is not asked here: \(\iota(U)\) is dense in \(V\), since for \(H = \widehat{(h_1,h_2,\dots)} \in V\) the estimate in the second half of (c) shows \(d_V(\iota(h_m), H) \to 0\). This makes \((V, d_V)\) the completion of \((U,d)\).

Exercises 6B

Problem (6B.1)

Show that if \(a, b \in \mathbb{R}\) with \(a + bi \neq 0\), then

\begin{equation*} \frac{1}{a + bi} = \frac{a}{a^2 + b^2} - \frac{b}{a^2 + b^2}\, i. \end{equation*}

Solution

Because \(a + bi \neq 0\), we do not have \(a = b = 0\); hence \(a^2 + b^2 > 0\) and the right side of the displayed equation makes sense.

Set

\begin{equation*} w = \frac{a}{a^2 + b^2} - \frac{b}{a^2 + b^2}\, i . \end{equation*}

To prove that \(w = \frac{1}{a+bi}\) it suffices, by the definition of the multiplicative inverse in the field \(\mathbb{C}\), to show that \((a + bi)w = 1\).

Using the definition of complex multiplication (6.17), which states that \((a + bi)(c + di) = (ac - bd) + (ad + bc)i\), with \(c = \frac{a}{a^2+b^2}\) and \(d = -\frac{b}{a^2+b^2}\), we get

\begin{equation*} \begin{aligned} (a + bi)w &= \left( a \cdot \frac{a}{a^2 + b^2} - b \cdot \left( -\frac{b}{a^2 + b^2} \right) \right)

  • \left( a \cdot \left( -\frac{b}{a^2 + b^2} \right) + b \cdot \frac{a}{a^2 + b^2} \right) i \\ &= \frac{a^2 + b^2}{a^2 + b^2} + \frac{-ab + ab}{a^2 + b^2}\, i \\ &= 1 + 0i \\ &= 1 . \end{aligned} \end{equation*}

Because \(\mathbb{C}\) is a field, the multiplicative inverse of the nonzero element \(a + bi\) is unique, and thus

\begin{equation*} \frac{1}{a + bi} = w = \frac{a}{a^2 + b^2} - \frac{b}{a^2 + b^2}\, i , \end{equation*}

as desired.

Equivalently, one can obtain the formula rather than merely verify it: by 6.24 we have \(z\bar{z} = |z|^2\) for every \(z \in \mathbb{C}\). Taking \(z = a + bi\), so that \(\bar{z} = a - bi\) and \(|z|^2 = a^2 + b^2 \neq 0\), dividing the equation \(z \bar z = |z|^2\) by \(z|z|^2\) gives

\begin{equation*} \frac{1}{a + bi} = \frac{\overline{a + bi}}{|a + bi|^2} = \frac{a - bi}{a^2 + b^2} = \frac{a}{a^2 + b^2} - \frac{b}{a^2 + b^2}\, i . \end{equation*}

Problem (6B.2)

Suppose \(z \in \mathbb{C}\). Prove that

\begin{equation*} \max\{|\operatorname{Re} z|, |\operatorname{Im} z|\} \leq |z| \leq \sqrt{2} \max\{|\operatorname{Re} z|, |\operatorname{Im} z|\}. \end{equation*}

Solution

Write \(z = a + bi\) with \(a, b \in \mathbb{R}\), so that \(\operatorname{Re} z = a\), \(\operatorname{Im} z = b\), and \(|z| = \sqrt{a^2 + b^2}\) by 6.18. Let

\begin{equation*} M = \max\{|a|, |b|\} . \end{equation*}

First we prove the left inequality. Because \(a^2 \geq 0\) and \(b^2 \geq 0\), we have

\begin{equation*} M^2 = \max\{a^2, b^2\} \leq a^2 + b^2 = |z|^2 . \end{equation*}

Here the first equality holds because \(t \mapsto t^2\) is increasing on \([0, \infty)\), so the larger of \(|a|, |b|\) has the larger square. Because \(M \geq 0\) and \(|z| \geq 0\), taking square roots (which preserves the order on \([0,\infty)\)) gives

\begin{equation*} \max\{|\operatorname{Re} z|, |\operatorname{Im} z|\} = M \leq |z| . \end{equation*}

Now we prove the right inequality. Because \(a^2 \leq M^2\) and \(b^2 \leq M^2\), we have

\begin{equation*} |z|^2 = a^2 + b^2 \leq 2M^2 = \left( \sqrt{2}\, M \right)^2 . \end{equation*}

Both \(|z|\) and \(\sqrt{2}M\) are nonnegative, so taking square roots gives

\begin{equation*} |z| \leq \sqrt{2} \, M = \sqrt{2} \max\{|\operatorname{Re} z|, |\operatorname{Im} z|\} , \end{equation*}

completing the proof.

Both inequalities are sharp: the left inequality becomes an equality when \(z\) is real or purely imaginary (for instance \(z = 1\)), and the right inequality becomes an equality when \(|\operatorname{Re} z| = |\operatorname{Im} z|\) (for instance \(z = 1 + i\), where \(|z| = \sqrt{2}\) and \(M = 1\)).

Problem (6B.3)

Suppose \(z \in \mathbb{C}\). Prove that

\begin{equation*} \frac{|\operatorname{Re} z| + |\operatorname{Im} z|}{\sqrt{2}} \leq |z| \leq |\operatorname{Re} z| + |\operatorname{Im} z|. \end{equation*}

Solution

Write \(z = a + bi\) with \(a, b \in \mathbb{R}\), so \(\operatorname{Re} z = a\), \(\operatorname{Im} z = b\), and \(|z| = \sqrt{a^2 + b^2}\) by 6.18.

For the right inequality: Because \(|a| \, |b| \geq 0\), we have

\begin{equation*} \left( |a| + |b| \right)^2 = a^2 + 2|a| \, |b| + b^2 \geq a^2 + b^2 = |z|^2 . \end{equation*}

Both \(|a| + |b|\) and \(|z|\) are nonnegative, so taking square roots gives \(|z| \leq |a| + |b|\), which is the right inequality.

For the left inequality: For all real numbers \(s, t\) we have \((|s| - |t|)^2 \geq 0\), which expands to \(2|s| \, |t| \leq s^2 + t^2\). Applying this with \(s = a\) and \(t = b\) gives

\begin{equation*} \left( |a| + |b| \right)^2 = a^2 + 2|a| \, |b| + b^2 \leq 2\left( a^2 + b^2 \right) = 2|z|^2 = \left( \sqrt{2}\, |z| \right)^2 . \end{equation*}

Again both quantities being squared are nonnegative, so taking square roots gives \(|a| + |b| \leq \sqrt{2}\,|z|\), that is,

\begin{equation*} \frac{|\operatorname{Re} z| + |\operatorname{Im} z|}{\sqrt{2}} \leq |z| , \end{equation*}

completing the proof.

Both inequalities are sharp. Equality holds on the right exactly when \(|a| \, |b| = 0\), that is, exactly when \(z\) is real or purely imaginary. Equality holds on the left exactly when \((|a| - |b|)^2 = 0\), that is, exactly when \(|\operatorname{Re} z| = |\operatorname{Im} z|\).

Problem (6B.4)

Suppose \(w, z \in \mathbb{C}\). Prove that \(|wz| = |w| \, |z|\) and \(|w + z| \leq |w| + |z|\).

Solution

Write \(w = a + bi\) and \(z = c + di\), where \(a, b, c, d \in \mathbb{R}\).

To prove multiplicativity of the absolute value, recall from 6.17 that

\begin{equation*} wz = (ac - bd) + (ad + bc) i . \end{equation*}

Hence, using the formula \(|u + vi|^2 = u^2 + v^2\) from 6.18 three times,

\begin{equation*} \begin{aligned} |wz|^2 &= (ac - bd)^2 + (ad + bc)^2 \\ &= a^2 c^2 - 2abcd + b^2 d^2 + a^2 d^2 + 2abcd + b^2 c^2 \\ &= a^2 c^2 + b^2 d^2 + a^2 d^2 + b^2 c^2 \\ &= \left( a^2 + b^2 \right) \left( c^2 + d^2 \right) \\ &= |w|^2 |z|^2 = \left( |w| \, |z| \right)^2 . \end{aligned} \end{equation*}

Because \(|wz|\) and \(|w| \, |z|\) are both nonnegative, taking square roots gives \(|wz| = |w| \, |z|\).

Now we prove the triangle inequality. We first record two facts.

First, if \(u \in \mathbb{C}\) then \(\operatorname{Re} u \leq |u|\). Indeed, \(\operatorname{Re} u \leq |\operatorname{Re} u| \leq |u|\), where the last inequality is the left inequality of Exercise 2 in this section (or directly: \((\operatorname{Re} u)^2 \leq (\operatorname{Re} u)^2 + (\operatorname{Im} u)^2 = |u|^2\)).

Second, \(|\bar{u}| = |u|\) for every \(u \in \mathbb{C}\) (this is one of the items of 6.24; it is immediate from 6.18 because \(\operatorname{Re}\bar u = \operatorname{Re} u\) and \(\operatorname{Im} \bar u = -\operatorname{Im} u\)).

By 6.24 we have \(u \bar{u} = |u|^2\) for all \(u \in \mathbb{C}\), and complex conjugation is additive and multiplicative. Thus

\begin{equation*} \begin{aligned} |w + z|^2 &= (w + z)\overline{(w + z)} \\ &= (w + z)(\bar{w} + \bar{z}) \\ &= w\bar{w} + z\bar{z} + w \bar{z} + \bar{w} z \\ &= |w|^2 + |z|^2 + w\bar{z} + \overline{w \bar{z}} , \end{aligned} \end{equation*}

where the last equality uses \(\overline{w\bar z} = \bar w \bar{\bar z} = \bar w z\) (multiplicativity of conjugation together with \(\bar{\bar z} = z\), both in 6.24).

By 6.24 we have \(u + \bar{u} = 2 \operatorname{Re} u\). Applying this with \(u = w \bar{z}\), and then using the two recorded facts along with the multiplicativity of the absolute value proved above, we obtain

\begin{equation*} w \bar{z} + \overline{w \bar{z}} = 2 \operatorname{Re}(w \bar{z}) \leq 2 |w \bar{z}| = 2 |w| \, |\bar{z}| = 2|w| \, |z| . \end{equation*}

Combining the two displays gives

\begin{equation*} |w + z|^2 \leq |w|^2 + 2|w| \, |z| + |z|^2 = \left( |w| + |z| \right)^2 . \end{equation*}

Because \(|w + z|\) and \(|w| + |z|\) are nonnegative, taking square roots yields

\begin{equation*} |w + z| \leq |w| + |z| , \end{equation*}

as desired.

Problem (6B.5)

Suppose \((X, \mathcal{S})\) is a measurable space and \(f \colon X \to \mathbb{C}\) is a complex-valued function. For conditions (b) and (c) below, identify \(\mathbb{C}\) with \(\mathbb{R}^2\). Prove that the following are equivalent.

(a)
\(f\) is \(\mathcal{S}\)-measurable.
(b)
\(f^{-1}(G) \in \mathcal{S}\) for every open set \(G\) in \(\mathbb{R}^2\).
(c)
\(f^{-1}(B) \in \mathcal{S}\) for every Borel set \(B \in \mathcal{B}_2\).
Solution

Throughout, \(\mathbb{C}\) is identified with \(\mathbb{R}^2\) by the bijection \(z \mapsto (\operatorname{Re} z, \operatorname{Im} z)\). Write

\begin{equation*} u = \operatorname{Re} f \quad\text{and}\quad v = \operatorname{Im} f , \end{equation*}

so that \(u, v \colon X \to \mathbb{R}\) and, under the identification, \(f(x) = \bigl( u(x), v(x) \bigr)\) for each \(x \in X\). Recall from 6.19 that \(f\) is \(\mathcal{S}\)-measurable means precisely that \(u\) and \(v\) are both \(\mathcal{S}\)-measurable real-valued functions.

We prove (a) \(\Rightarrow\) (c) \(\Rightarrow\) (b) \(\Rightarrow\) (a).

First suppose (a) holds, so \(u\) and \(v\) are \(\mathcal{S}\)-measurable. Let

\begin{equation*} \mathcal{T} = \left\{ B \subseteq \mathbb{R}^2 : f^{-1}(B) \in \mathcal{S} \right\} . \end{equation*}

We claim \(\mathcal{T}\) is a \(\sigma\)-algebra on \(\mathbb{R}^2\). Indeed \(f^{-1}(\varnothing) = \varnothing \in \mathcal{S}\), so \(\varnothing \in \mathcal{T}\). If \(B \in \mathcal{T}\), then by 2.33(a) we have \(f^{-1}(\mathbb{R}^2 \setminus B) = X \setminus f^{-1}(B) \in \mathcal{S}\), so \(\mathbb{R}^2 \setminus B \in \mathcal{T}\); thus \(\mathcal{T}\) is closed under complementation. If \(B_1, B_2, \ldots \in \mathcal{T}\), then by 2.33(b) we have \(f^{-1}\bigl( \bigcup_{k=1}^{\infty} B_k \bigr) = \bigcup_{k=1}^{\infty} f^{-1}(B_k) \in \mathcal{S}\), so \(\mathcal{T}\) is closed under countable unions. Hence \(\mathcal{T}\) is a \(\sigma\)-algebra on \(\mathbb{R}^2\).

Next we show that \(\mathcal{T}\) contains every open cube in \(\mathbb{R}^2\). Recall that for \(x \in \mathbb{R}^2\) and \(\delta > 0\) the open cube \(B(x, \delta)\) consists of the points \(y\) with \(\| y - x \|_{\infty} < \delta\); thus an open cube in \(\mathbb{R}^2\) is exactly a set of the form \(C = I_1 \times I_2\), where \(I_1\) and \(I_2\) are bounded open intervals of the same length (namely \(2\delta\)). For such \(C\),

\begin{equation*} f^{-1}( C) = \left\{ x \in X : u(x) \in I_1 \text{ and } v(x) \in I_2 \right\} = u^{-1}(I_1) \cap v^{-1}(I_2) . \end{equation*}

Because \(I_1\) and \(I_2\) are open subsets of \(\mathbb{R}\), they are Borel sets, and thus \(u^{-1}(I_1) \in \mathcal{S}\) and \(v^{-1}(I_2) \in \mathcal{S}\) by the definition (2.35) of \(\mathcal{S}\)-measurability of the real-valued functions \(u\) and \(v\). Because \(\mathcal{S}\) is a \(\sigma\)-algebra it is closed under finite intersections (2.25(b)), so \(f^{-1}( C) \in \mathcal{S}\), that is, \(C \in \mathcal{T}\).

Thus \(\mathcal{T}\) is a \(\sigma\)-algebra on \(\mathbb{R}^2\) containing all open cubes in \(\mathbb{R}^2\). By 5.38(b), \(\mathcal{B}_2\) is the smallest such \(\sigma\)-algebra, and therefore \(\mathcal{B}_2 \subseteq \mathcal{T}\). In other words, \(f^{-1}(B) \in \mathcal{S}\) for every \(B \in \mathcal{B}_2\), which is (c).

Now suppose (c) holds. By the definition of \(\mathcal{B}_2\) (see 5.37), every open subset \(G\) of \(\mathbb{R}^2\) belongs to \(\mathcal{B}_2\). Hence \(f^{-1}(G) \in \mathcal{S}\) for every open \(G \subseteq \mathbb{R}^2\), which is (b).

Finally, suppose (b) holds. Fix \(a \in \mathbb{R}\). The set \((a, \infty) \times \mathbb{R}\) is an open subset of \(\mathbb{R}^2\), and

\begin{equation*} u^{-1}\bigl( (a, \infty) \bigr) = \left\{ x \in X : u(x) > a \right\} = f^{-1}\bigl( (a, \infty) \times \mathbb{R} \bigr) , \end{equation*}

which is in \(\mathcal{S}\) by (b). Because \(a \in \mathbb{R}\) was arbitrary, 2.39 implies that \(u\) is an \(\mathcal{S}\)-measurable function. Similarly, \(\mathbb{R} \times (a, \infty)\) is open in \(\mathbb{R}^2\) and

\begin{equation*} v^{-1}\bigl( (a, \infty) \bigr) = f^{-1}\bigl( \mathbb{R} \times (a, \infty) \bigr) \in \mathcal{S} \end{equation*}

for all \(a \in \mathbb{R}\), so 2.39 implies that \(v\) is \(\mathcal{S}\)-measurable. Because \(u = \operatorname{Re} f\) and \(v = \operatorname{Im} f\) are both \(\mathcal{S}\)-measurable, 6.19 tells us that \(f\) is \(\mathcal{S}\)-measurable, which is (a).

Having proved (a) \(\Rightarrow\) (c) \(\Rightarrow\) (b) \(\Rightarrow\) (a), we conclude that the three conditions are equivalent.

Problem (6B.6)

Suppose \((X, \mathcal{S})\) is a measurable space and \(f, g \colon X \to \mathbb{C}\) are \(\mathcal{S}\)-measurable. Prove that

(a)
\(f + g\), \(f - g\), and \(fg\) are \(\mathcal{S}\)-measurable functions;
(b)
if \(g(x) \neq 0\) for all \(x \in X\), then \(\frac{f}{g}\) is an \(\mathcal{S}\)-measurable function.
Solution

Write

\begin{equation*} f = u_1 + i v_1 \quad\text{and}\quad g = u_2 + i v_2 , \end{equation*}

where \(u_1 = \operatorname{Re} f\), \(v_1 = \operatorname{Im} f\), \(u_2 = \operatorname{Re} g\), \(v_2 = \operatorname{Im} g\) are real-valued functions on \(X\). By 6.19, the hypothesis that \(f\) and \(g\) are \(\mathcal{S}\)-measurable means exactly that the four real-valued functions \(u_1, v_1, u_2, v_2\) are \(\mathcal{S}\)-measurable. Also by 6.19, to prove that a complex-valued function on \(X\) is \(\mathcal{S}\)-measurable it suffices to prove that its real and imaginary parts are \(\mathcal{S}\)-measurable real-valued functions. Throughout we use 2.46, which states that sums, differences, products, and (where the denominator never vanishes) quotients of \(\mathcal{S}\)-measurable real-valued functions are \(\mathcal{S}\)-measurable.

For (a), the definition of addition in \(\mathbb{C}\) (6.17) gives

\begin{equation*} \operatorname{Re}(f + g) = u_1 + u_2 , \qquad \operatorname{Im}(f + g) = v_1 + v_2 , \end{equation*}

where these equalities hold pointwise on \(X\). Both right sides are \(\mathcal{S}\)-measurable by 2.46(a). Hence \(f + g\) is \(\mathcal{S}\)-measurable by 6.19. In the same way,

\begin{equation*} \operatorname{Re}(f - g) = u_1 - u_2 , \qquad \operatorname{Im}(f - g) = v_1 - v_2 , \end{equation*}

which are \(\mathcal{S}\)-measurable by 2.46(a); hence \(f - g\) is \(\mathcal{S}\)-measurable.

For the product, the definition of multiplication in \(\mathbb{C}\) (6.17) gives, pointwise on \(X\),

\begin{equation*} fg = (u_1 u_2 - v_1 v_2) + (u_1 v_2 + v_1 u_2) i , \end{equation*}

so

\begin{equation*} \operatorname{Re}(fg) = u_1 u_2 - v_1 v_2 , \qquad \operatorname{Im}(fg) = u_1 v_2 + v_1 u_2 . \end{equation*}

Each of \(u_1 u_2\), \(v_1 v_2\), \(u_1 v_2\), \(v_1 u_2\) is \(\mathcal{S}\)-measurable by 2.46(a) (products), and hence so are the displayed difference and sum, again by 2.46(a). Thus \(fg\) is \(\mathcal{S}\)-measurable by 6.19, completing the proof of (a).

For (b), suppose \(g(x) \neq 0\) for all \(x \in X\). Define \(h \colon X \to \mathbb{R}\) by

\begin{equation*} h = |g|^2 = u_2^{\,2} + v_2^{\,2} . \end{equation*}

Then \(u_2^{\,2}\) and \(v_2^{\,2}\) are \(\mathcal{S}\)-measurable by 2.45, and hence \(h\) is \(\mathcal{S}\)-measurable by 2.46(a). Moreover \(h(x) > 0\) for every \(x \in X\): if \(h(x) = 0\) then \(u_2(x) = v_2(x) = 0\), so \(g(x) = 0\), contradicting the hypothesis. In particular \(h(x) \neq 0\) for all \(x \in X\).

By 6.24 we have \(g \bar{g} = |g|^2 = h\), and \(h\) is nowhere zero, so dividing by \(gh\) shows that pointwise on \(X\)

\begin{equation*} \frac{f}{g} = \frac{f \bar{g}}{g \bar{g}} = \frac{f \bar{g}}{h} . \end{equation*}

By 6.23, \(\bar{g} = u_2 - i v_2\), so by the multiplication rule in 6.17,

\begin{equation*} f \bar{g} = (u_1 + i v_1)(u_2 - i v_2) = (u_1 u_2 + v_1 v_2) + (v_1 u_2 - u_1 v_2) i . \end{equation*}

Therefore

\begin{equation*} \operatorname{Re}\frac{f}{g} = \frac{u_1 u_2 + v_1 v_2}{u_2^{\,2} + v_2^{\,2}} , \qquad \operatorname{Im}\frac{f}{g} = \frac{v_1 u_2 - u_1 v_2}{u_2^{\,2} + v_2^{\,2}} , \end{equation*}

where we have used that dividing a complex number by the nonzero real number \(h\) divides its real part and its imaginary part by \(h\).

The two numerators above are \(\mathcal{S}\)-measurable by 2.46(a), and the denominator \(h\) is an \(\mathcal{S}\)-measurable function that is never \(0\). Hence both quotients are \(\mathcal{S}\)-measurable by 2.46(b). Thus \(\operatorname{Re}\frac{f}{g}\) and \(\operatorname{Im}\frac{f}{g}\) are \(\mathcal{S}\)-measurable, and so \(\frac{f}{g}\) is an \(\mathcal{S}\)-measurable function by 6.19.

Problem (6B.7)

Suppose \((X, \mathcal{S})\) is a measurable space and \(f_1, f_2, \ldots\) is a sequence of \(\mathcal{S}\)-measurable functions from \(X\) to \(\mathbb{C}\). Suppose \(\lim_{k \to \infty} f_k(x)\) exists for each \(x \in X\). Define \(f \colon X \to \mathbb{C}\) by

\begin{equation*} f(x) = \lim_{k \to \infty} f_k(x) . \end{equation*}

Prove that \(f\) is an \(\mathcal{S}\)-measurable function.

Solution

The key point is that convergence in \(\mathbb{C}\) is equivalent to convergence of the real parts together with convergence of the imaginary parts. Section 6B records this after 6.18: if \(z_1, z_2, \ldots \in \mathbb{C}\) and \(L \in \mathbb{C}\), then

\begin{equation*} \lim_{k \to \infty} z_k = L \iff \lim_{k \to \infty} \operatorname{Re} z_k = \operatorname{Re} L \ \text{ and } \ \lim_{k \to \infty} \operatorname{Im} z_k = \operatorname{Im} L . \end{equation*}

For completeness, here is the proof of the direction we need. Suppose \(\lim_{k \to \infty} z_k = L\), meaning \(\lim_{k \to \infty} |z_k - L| = 0\) (see 6.18). Because \(\operatorname{Re}\) and \(\operatorname{Im}\) are additive, \(\operatorname{Re} z_k - \operatorname{Re} L = \operatorname{Re}(z_k - L)\) and \(\operatorname{Im} z_k - \operatorname{Im} L = \operatorname{Im}(z_k - L)\). The left inequality of Exercise 2 in this section gives

\begin{equation*} \left| \operatorname{Re} z_k - \operatorname{Re} L \right| \leq |z_k - L| \quad\text{and}\quad \left| \operatorname{Im} z_k - \operatorname{Im} L \right| \leq |z_k - L| . \end{equation*}

Letting \(k \to \infty\) and using the Squeeze Theorem for real sequences, we get \(\lim_{k \to \infty} \operatorname{Re} z_k = \operatorname{Re} L\) and \(\lim_{k \to \infty} \operatorname{Im} z_k = \operatorname{Im} L\).

Now fix \(x \in X\). Applying the above with \(z_k = f_k(x)\) and \(L = f(x)\) (the limit exists by hypothesis) shows that

\begin{equation*} \lim_{k \to \infty} (\operatorname{Re} f_k)(x) = (\operatorname{Re} f)(x) \quad\text{and}\quad \lim_{k \to \infty} (\operatorname{Im} f_k)(x) = (\operatorname{Im} f)(x) . \end{equation*}

In particular, both of these limits exist for every \(x \in X\).

Each \(f_k\) is an \(\mathcal{S}\)-measurable complex-valued function, so by 6.19 the real-valued functions \(\operatorname{Re} f_k\) and \(\operatorname{Im} f_k\) are \(\mathcal{S}\)-measurable. Thus \(\operatorname{Re} f_1, \operatorname{Re} f_2, \ldots\) is a sequence of \(\mathcal{S}\)-measurable functions from \(X\) to \(\mathbb{R}\) whose pointwise limit exists at every point of \(X\) and equals \(\operatorname{Re} f\). Hence 2.48 implies that \(\operatorname{Re} f\) is an \(\mathcal{S}\)-measurable function. The same argument applied to \(\operatorname{Im} f_1, \operatorname{Im} f_2, \ldots\) shows that \(\operatorname{Im} f\) is an \(\mathcal{S}\)-measurable function.

Because \(\operatorname{Re} f\) and \(\operatorname{Im} f\) are both \(\mathcal{S}\)-measurable, 6.19 tells us that \(f\) is an \(\mathcal{S}\)-measurable function, as desired.

Problem (6B.8)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(f : X \to \mathbb{C}\) is an \(\mathcal{S}\)-measurable function such that \(\int |f| \, d\mu < \infty\). Prove that if \(\alpha \in \mathbb{C}\), then

\begin{equation*} \int \alpha f \, d\mu = \alpha \int f \, d\mu . \end{equation*}

Solution

Write \(u = \operatorname{Re} f\) and \(v = \operatorname{Im} f\), so that \(u, v : X \to \mathbb{R}\) are \(\mathcal{S}\)-measurable (by 6.19) and \(f = u + iv\). Write \(\alpha = a + bi\) with \(a, b \in \mathbb{R}\).

First we check that every integral appearing below is defined and finite. For each \(x \in X\) we have \(|u(x)| \le |f(x)|\) and \(|v(x)| \le |f(x)|\), because \(|f(x)| = \sqrt{u(x)^2 + v(x)^2}\). The functions \(|u|, |v|, |f|\) are nonnegative and \(\mathcal{S}\)-measurable (measurability of \(|f|\) is the case \(p = 1\) of 6.20), so 3.8 gives

\begin{equation*} \int |u| \, d\mu \le \int |f| \, d\mu < \infty \quad\text{and}\quad \int |v| \, d\mu \le \int |f| \, d\mu < \infty . \end{equation*}

Now \(u^{+} \le |u|\) and \(u^{-} \le |u|\) pointwise, so 3.8 gives \(\int u^{+} \, d\mu < \infty\) and \(\int u^{-} \, d\mu < \infty\); hence \(\int u \, d\mu = \int u^{+} \, d\mu - \int u^{-} \, d\mu\) is defined and is a real number. The same argument applies to \(v\). Thus \(\int f \, d\mu\) is defined by 6.21 as

\begin{equation*} \int f \, d\mu = \int u \, d\mu + i \int v \, d\mu . \end{equation*}

Next we check that \(\int \alpha f \, d\mu\) is defined. Multiplying out with the definition 6.17 of complex multiplication,

\begin{equation*} \alpha f = (a + bi)(u + iv) = (au - bv) + i\,(av + bu), \end{equation*}

so

\begin{equation*} \operatorname{Re}(\alpha f) = au - bv \quad\text{and}\quad \operatorname{Im}(\alpha f) = av + bu . \end{equation*}

Constant functions are \(\mathcal{S}\)-measurable, so \(au\), \(bv\), \(av\), \(bu\) are \(\mathcal{S}\)-measurable by 2.46(a), and therefore \(au - bv\) and \(av + bu\) are \(\mathcal{S}\)-measurable, again by 2.46(a). Thus \(\alpha f\) is an \(\mathcal{S}\)-measurable function from \(X\) to \(\mathbb{C}\), by 6.19. Moreover \(|\alpha f(x)| = |\alpha| \, |f(x)|\) for every \(x \in X\), by the multiplicativity of the absolute value on \(\mathbb{C}\) (Exercise 4 in this section; it also follows from a direct computation with 6.17 and 6.18). Since \(|f|\) is nonnegative and \(\mathcal{S}\)-measurable, its integral is defined, so 3.20 with \(c = |\alpha|\) gives

\begin{equation*} \int |\alpha f| \, d\mu = |\alpha| \int |f| \, d\mu < \infty . \end{equation*}

Hence \(\int \alpha f \, d\mu\) is defined, and 6.21 gives

\begin{equation*} \int \alpha f \, d\mu = \int (au - bv) \, d\mu + i \int (av + bu) \, d\mu . \end{equation*}

Now we evaluate the two real integrals on the right. Each of \(au\), \(bv\), \(av\), \(bu\) has absolute value with a finite integral: for instance \(\int |au| \, d\mu = |a| \int |u| \, d\mu < \infty\) by 3.20. Hence additivity of integration (3.21) applies to each of the two sums, and homogeneity of integration (3.20) applies to each real scalar multiple. This gives

\begin{equation*} \begin{aligned} \int (au - bv) \, d\mu &= \int au \, d\mu + \int (-b)v \, d\mu = a \int u \, d\mu - b \int v \, d\mu, \\ \int (av + bu) \, d\mu &= \int av \, d\mu + \int bu \, d\mu = a \int v \, d\mu + b \int u \, d\mu . \end{aligned} \end{equation*}

Therefore

\begin{equation*} \int \alpha f \, d\mu = \left( a \int u \, d\mu - b \int v \, d\mu \right) + i \left( a \int v \, d\mu + b \int u \, d\mu \right). \end{equation*}

On the other hand, using 6.21 and the definition 6.17 of complex multiplication,

\begin{equation*} \alpha \int f \, d\mu = (a + bi)\left( \int u \, d\mu + i \int v \, d\mu \right) = \left( a \int u \, d\mu - b \int v \, d\mu \right) + i \left( a \int v \, d\mu + b \int u \, d\mu \right). \end{equation*}

The two displayed expressions are identical, so \(\int \alpha f \, d\mu = \alpha \int f \, d\mu\), as desired.

Problem (6B.9)

Suppose \(V\) is a vector space. Show that the intersection of every collection of subspaces of \(V\) is a subspace of \(V\).

Solution

Let \(\mathcal{A}\) be a collection of subspaces of \(V\) and let

\begin{equation*} U = \bigcap_{U^{\prime} \in \mathcal{A}} U^{\prime} \end{equation*}

denote the intersection of the collection, the intersection being taken inside \(V\) (with the usual convention that the intersection of the empty collection of subsets of \(V\) is \(V\) itself, which is a subspace of \(V\); the argument below covers that degenerate case as well, since each of its three verifications is then vacuously about \(V\)).

Note first that \(U \subseteq V\): this is immediate when \(\mathcal{A} \neq \emptyset\), since \(U \subseteq U^{\prime} \subseteq V\) for any \(U^{\prime} \in \mathcal{A}\), and it holds by the stated convention when \(\mathcal{A} = \emptyset\). We now verify the three conditions of 6.31.

Additive identity. Each \(U^{\prime} \in \mathcal{A}\) is a subspace of \(V\), so the forward direction of 6.31 gives \(0 \in U^{\prime}\), where \(0\) denotes the additive identity of \(V\). Since \(0 \in U^{\prime}\) for every \(U^{\prime} \in \mathcal{A}\), we have \(0 \in U\). In particular \(U\) is nonempty.

Closed under addition. Suppose \(f, g \in U\). Let \(U^{\prime} \in \mathcal{A}\). Then \(f \in U^{\prime}\) and \(g \in U^{\prime}\) by the definition of intersection, and \(U^{\prime}\) is closed under addition by 6.31, so \(f + g \in U^{\prime}\). As \(U^{\prime} \in \mathcal{A}\) was arbitrary, \(f + g\) belongs to every member of \(\mathcal{A}\), which means \(f + g \in U\).

Closed under scalar multiplication. Suppose \(\alpha \in \mathbf{F}\) and \(f \in U\). Let \(U^{\prime} \in \mathcal{A}\). Then \(f \in U^{\prime}\), and \(U^{\prime}\) is closed under scalar multiplication by 6.31, so \(\alpha f \in U^{\prime}\). As \(U^{\prime} \in \mathcal{A}\) was arbitrary, \(\alpha f \in U\).

Thus \(U\) is a subset of \(V\) satisfying the three conditions of 6.31, and therefore the reverse direction of 6.31 shows that \(U\) is a subspace of \(V\).

Problem (6B.10)

Suppose \(V\) and \(W\) are vector spaces. Define \(V \times W\) by

\begin{equation*} V \times W = \{ (f, g) : f \in V \text{ and } g \in W \}. \end{equation*}

Define addition and scalar multiplication on \(V \times W\) by

\begin{equation*} (f_1, g_1) + (f_2, g_2) = (f_1 + f_2,\ g_1 + g_2) \quad\text{and}\quad \alpha(f, g) = (\alpha f,\ \alpha g). \end{equation*}

Prove that \(V \times W\) is a vector space with these operations.

Solution

Here \(V\) and \(W\) are vector spaces over the same field \(\mathbf{F}\) (recall 6.25: \(\mathbf{F}\) denotes \(\mathbb{R}\) or \(\mathbb{C}\)), which is what makes the definition of \(\alpha(f,g)\) meaningful. Throughout we use repeatedly that two ordered pairs are equal if and only if their first coordinates are equal and their second coordinates are equal.

The operations are well defined. If \((f_1, g_1), (f_2, g_2) \in V \times W\), then \(f_1 + f_2 \in V\) because addition on \(V\) maps into \(V\), and \(g_1 + g_2 \in W\) likewise; hence \((f_1 + f_2, g_1 + g_2) \in V \times W\). Similarly, if \(\alpha \in \mathbf{F}\) and \((f,g) \in V \times W\), then \(\alpha f \in V\) and \(\alpha g \in W\), so \((\alpha f, \alpha g) \in V \times W\). Thus the two formulas do define an addition and a scalar multiplication on \(V \times W\) in the sense of 6.26.

We now verify each property in the definition 6.27 of a vector space. In every case the verification is carried out coordinatewise, using the corresponding property in \(V\) for the first coordinate and in \(W\) for the second.

Commutativity. For \((f_1, g_1), (f_2, g_2) \in V \times W\),

\begin{equation*} (f_1, g_1) + (f_2, g_2) = (f_1 + f_2,\ g_1 + g_2) = (f_2 + f_1,\ g_2 + g_1) = (f_2, g_2) + (f_1, g_1), \end{equation*}

where the middle equality uses commutativity of addition in \(V\) and in \(W\).

Associativity of addition. For \((f_1, g_1), (f_2, g_2), (f_3, g_3) \in V \times W\),

\begin{equation*} \begin{aligned} \bigl( (f_1, g_1) + (f_2, g_2) \bigr) + (f_3, g_3) &= (f_1 + f_2,\ g_1 + g_2) + (f_3, g_3) \\ &= \bigl( (f_1 + f_2) + f_3,\ (g_1 + g_2) + g_3 \bigr) \\ &= \bigl( f_1 + (f_2 + f_3),\ g_1 + (g_2 + g_3) \bigr) \\ &= (f_1, g_1) + (f_2 + f_3,\ g_2 + g_3) \\ &= (f_1, g_1) + \bigl( (f_2, g_2) + (f_3, g_3) \bigr), \end{aligned} \end{equation*}

where the third equality uses associativity of addition in \(V\) and in \(W\).

Associativity of scalar multiplication. For \(\alpha, \beta \in \mathbf{F}\) and \((f,g) \in V \times W\),

\begin{equation*} (\alpha\beta)(f, g) = \bigl( (\alpha\beta) f,\ (\alpha\beta) g \bigr) = \bigl( \alpha(\beta f),\ \alpha(\beta g) \bigr) = \alpha(\beta f,\ \beta g) = \alpha \bigl( \beta (f,g) \bigr), \end{equation*}

where the second equality uses the corresponding associativity property in \(V\) and in \(W\).

Additive identity. Let \(0_V\) and \(0_W\) denote the additive identities of \(V\) and \(W\), and set \(0 = (0_V, 0_W) \in V \times W\). For every \((f,g) \in V \times W\),

\begin{equation*} (f, g) + (0_V, 0_W) = (f + 0_V,\ g + 0_W) = (f, g). \end{equation*}

Thus \((0_V, 0_W)\) is an additive identity for \(V \times W\).

Additive inverse. Suppose \((f, g) \in V \times W\). Because \(V\) is a vector space, there exists \(f^{\prime} \in V\) with \(f + f^{\prime} = 0_V\); because \(W\) is a vector space, there exists \(g^{\prime} \in W\) with \(g + g^{\prime} = 0_W\). Then \((f^{\prime}, g^{\prime}) \in V \times W\) and

\begin{equation*} (f, g) + (f^{\prime}, g^{\prime}) = (f + f^{\prime},\ g + g^{\prime}) = (0_V, 0_W). \end{equation*}

Thus every element of \(V \times W\) has an additive inverse in \(V \times W\).

Multiplicative identity. For \((f, g) \in V \times W\),

\begin{equation*} 1(f, g) = (1f,\ 1g) = (f, g), \end{equation*}

using the multiplicative identity property in \(V\) and in \(W\).

Distributive properties. For \(\alpha \in \mathbf{F}\) and \((f_1, g_1), (f_2, g_2) \in V \times W\),

\begin{equation*} \begin{aligned} \alpha \bigl( (f_1, g_1) + (f_2, g_2) \bigr) &= \alpha (f_1 + f_2,\ g_1 + g_2) \\ &= \bigl( \alpha(f_1 + f_2),\ \alpha(g_1 + g_2) \bigr) \\ &= (\alpha f_1 + \alpha f_2,\ \alpha g_1 + \alpha g_2) \\ &= (\alpha f_1,\ \alpha g_1) + (\alpha f_2,\ \alpha g_2) \\ &= \alpha(f_1, g_1) + \alpha(f_2, g_2), \end{aligned} \end{equation*}

where the third equality uses the first distributive property in \(V\) and in \(W\). Also, for \(\alpha, \beta \in \mathbf{F}\) and \((f, g) \in V \times W\),

\begin{equation*} \begin{aligned} (\alpha + \beta)(f, g) &= \bigl( (\alpha + \beta) f,\ (\alpha + \beta) g \bigr) \\ &= (\alpha f + \beta f,\ \alpha g + \beta g) \\ &= (\alpha f,\ \alpha g) + (\beta f,\ \beta g) \\ &= \alpha(f, g) + \beta(f, g), \end{aligned} \end{equation*}

where the second equality uses the second distributive property in \(V\) and in \(W\).

All the properties required by 6.27 hold, so \(V \times W\), with the addition and scalar multiplication defined above, is a vector space over \(\mathbf{F}\).

Exercises 6C

Problem (6C.1)

Show that the map \(f \mapsto \|f\|\) from a normed vector space \(V\) to \(\mathbf{F}\) is continuous (where the norm on \(\mathbf{F}\) is the usual absolute value).

Solution

Let \((V, \|\cdot\|)\) be a normed vector space and let \(N \colon V \to \mathbf{F}\) be defined by \(N(f) = \|f\|\). Recall from 6.36 that the metric on \(V\) is \(d(f,g) = \|f - g\|\), and the metric on \(\mathbf{F}\) is \(d(s,t) = |s - t|\).

First we prove the reverse triangle inequality: if \(f, g \in V\), then

\begin{equation*} \bigl| \|f\| - \|g\| \bigr| \le \|f - g\|. \end{equation*}

To see this, note that the triangle inequality (see 6.33) gives

\begin{equation*} \|f\| = \|(f - g) + g\| \le \|f - g\| + \|g\|, \end{equation*}

and hence \(\|f\| - \|g\| \le \|f - g\|\). Interchanging the roles of \(f\) and \(g\) gives \(\|g\| - \|f\| \le \|g - f\|\). By homogeneity (see 6.33) with \(\alpha = -1\) we have

\begin{equation*} \|g - f\| = \|(-1)(f - g)\| = |-1| \, \|f - g\| = \|f - g\|. \end{equation*}

Thus \(\|g\| - \|f\| \le \|f - g\|\) as well. Because \(\bigl| \|f\| - \|g\| \bigr|\) equals either \(\|f\| - \|g\|\) or \(\|g\| - \|f\|\), the displayed inequality holds.

Now we verify continuity at an arbitrary \(f \in V\), using the definition of continuity for metric spaces (see 6.10). Suppose \(\varepsilon > 0\). Take \(\delta = \varepsilon\). If \(g \in V\) satisfies \(\|g - f\| < \delta\), then the reverse triangle inequality gives

\begin{equation*} \bigl| N(g) - N(f) \bigr| = \bigl| \|g\| - \|f\| \bigr| \le \|g - f\| < \delta = \varepsilon. \end{equation*}

Thus \(N\) is continuous at \(f\). Because \(f \in V\) was arbitrary, \(N\) is continuous on \(V\).

In fact the estimate above shows more: \(N\) is Lipschitz with Lipschitz constant \(1\), and hence uniformly continuous on \(V\).

Note that \(N\) is not a linear map (indeed \(N\) takes only nonnegative real values, so it is not linear unless \(V = \{0\}\)), so 6.48 does not apply here; the direct argument above is needed.

Problem (6C.2)

Prove that if \(V\) is a normed vector space, \(f \in V\), and \(r > 0\), then

\begin{equation*} \overline{B(f,r)} = \overline{B}(f,r). \end{equation*}

Solution

Here \(B(f,r) = \{g \in V : \|f - g\| < r\}\) is the open ball and \(\overline{B}(f,r) = \{g \in V : \|f - g\| \le r\}\) is the closed ball (see 6.3, interpreted via the metric \(d(g,h) = \|g - h\|\) of 6.36), while \(\overline{B(f,r)}\) denotes the closure of the open ball (see 6.7). We prove the two inclusions separately, using the characterization of closure by limits of sequences given in 6.9(a).

First suppose \(g \in \overline{B(f,r)}\). By 6.9(a) there is a sequence \(g_1, g_2, \ldots\) in \(B(f,r)\) with \(\lim_{k \to \infty} g_k = g\), which in the normed vector space context means \(\lim_{k \to \infty} \|g_k - g\| = 0\). For each \(k \in \mathbf{Z}^+\) the triangle inequality gives

\begin{equation*} \|f - g\| \le \|f - g_k\| + \|g_k - g\| < r + \|g_k - g\|. \end{equation*}

Letting \(k \to \infty\) in the inequality \(\|f - g\| < r + \|g_k - g\|\) shows that \(\|f - g\| \le r\). Thus \(g \in \overline{B}(f,r)\), proving that \(\overline{B(f,r)} \subseteq \overline{B}(f,r)\).

Now suppose \(g \in \overline{B}(f,r)\), so \(\|f - g\| \le r\). For \(k \in \mathbf{Z}^+\) define

\begin{equation*} g_k = f + \Bigl(1 - \tfrac{1}{k}\Bigr)(g - f). \end{equation*}

Because \(1 - \frac{1}{k} \in [0,1)\), homogeneity of the norm gives

\begin{equation*} \|g_k - f\| = \Bigl(1 - \tfrac{1}{k}\Bigr) \|g - f\| \le \Bigl(1 - \tfrac{1}{k}\Bigr) r < r, \end{equation*}

so \(g_k \in B(f,r)\) for each \(k \in \mathbf{Z}^+\). Furthermore

\begin{equation*} \|g_k - g\| = \Bigl\| \Bigl(1 - \tfrac{1}{k}\Bigr)(g - f) - (g - f) \Bigr\| = \tfrac{1}{k} \|g - f\| \le \tfrac{r}{k}, \end{equation*}

which tends to \(0\) as \(k \to \infty\). Hence \(\lim_{k \to \infty} g_k = g\) with all \(g_k \in B(f,r)\), and so \(g \in \overline{B(f,r)}\) by 6.9(a). This proves that \(\overline{B}(f,r) \subseteq \overline{B(f,r)}\).

The two inclusions together give \(\overline{B(f,r)} = \overline{B}(f,r)\).

Note that the vector space structure is essential here; the corresponding statement is false in a general metric space. For example, if \(V\) is a set with at least two elements equipped with the discrete metric (\(d(g,h) = 1\) when \(g \ne h\)), then \(B(f,1) = \{f\}\), which is closed, so \(\overline{B(f,1)} = \{f\}\), whereas \(\overline{B}(f,1) = V\). The second half of the proof above used scalar multiplication to slide points of the closed ball inward, which is exactly what the discrete metric space lacks.

Problem (6C.3)

Show that the functions defined in the last two bullet points of Example 6.35 are not norms.

Solution

Recall the two functions in question. Suppose \(n \in \mathbf{Z}^+\). The second-to-last bullet point of 6.35 defines \(\|\cdot\|\) on \(\mathbf{F}^n\) by

\begin{equation*} \|(a_1, \ldots, a_n)\| = |a_1|^{1/2} + \cdots + |a_n|^{1/2}, \end{equation*}

and the last bullet point defines \(\|\cdot\|_{1/2}\) on \(\mathbf{F}^n\) by

\begin{equation*} \|(a_1, \ldots, a_n)\|_{1/2} = \bigl( |a_1|^{1/2} + \cdots + |a_n|^{1/2} \bigr)^2 . \end{equation*}

By 6.33, to show that a function is not a norm it suffices to exhibit a failure of any one of the three defining conditions.

First consider \(\|\cdot\|\) as defined in the second-to-last bullet point. We show that homogeneity fails. Let \(f = (1, 0, \ldots, 0) \in \mathbf{F}^n\) and let \(\alpha = 4\). Then

\begin{equation*} \|f\| = |1|^{1/2} = 1 \qquad \text{and} \qquad \|\alpha f\| = \|(4,0,\ldots,0)\| = |4|^{1/2} = 2 . \end{equation*}

Thus \(\|\alpha f\| = 2\), while \(|\alpha| \, \|f\| = 4 \cdot 1 = 4\). Hence \(\|\alpha f\| \ne |\alpha| \, \|f\|\), so the homogeneity condition in 6.33 fails and \(\|\cdot\|\) is not a norm on \(\mathbf{F}^n\).

More generally, for every \(\alpha \in \mathbf{F}\) and every \((a_1, \ldots, a_n) \in \mathbf{F}^n\) we have

\begin{equation*} \|\alpha(a_1, \ldots, a_n)\| = \sum_{k=1}^{n} |\alpha a_k|^{1/2} = |\alpha|^{1/2} \sum_{k=1}^{n} |a_k|^{1/2} = |\alpha|^{1/2} \, \|(a_1, \ldots, a_n)\| , \end{equation*}

so this function scales like \(|\alpha|^{1/2}\) rather than like \(|\alpha|\); the two agree only when \(|\alpha| \in \{0, 1\}\) or when the vector is \(0\).

Now consider \(\|\cdot\|_{1/2}\) as defined in the last bullet point, and suppose \(n > 1\). (This hypothesis is necessary: when \(n = 1\) we have \(\|(a_1)\|_{1/2} = \bigl(|a_1|^{1/2}\bigr)^2 = |a_1|\), which is the usual absolute value norm on \(\mathbf{F}\).) Note first that \(\|\cdot\|_{1/2}\) does satisfy homogeneity, since

\begin{equation*} \|\alpha(a_1, \ldots, a_n)\|_{1/2} = \Bigl( |\alpha|^{1/2} \sum_{k=1}^{n} |a_k|^{1/2} \Bigr)^{2} = |\alpha| \, \|(a_1, \ldots, a_n)\|_{1/2}, \end{equation*}

and it is positive definite, because the sum \(\sum_{k=1}^n |a_k|^{1/2}\) is \(0\) if and only if every \(a_k\) is \(0\). What fails is the triangle inequality. Let

\begin{equation*} f = (1, 0, 0, \ldots, 0) \qquad \text{and} \qquad g = (0, 1, 0, \ldots, 0), \end{equation*}

both of which lie in \(\mathbf{F}^n\) because \(n \ge 2\). Then

\begin{equation*} \|f\|_{1/2} = \bigl(|1|^{1/2}\bigr)^2 = 1, \qquad \|g\|_{1/2} = 1, \end{equation*}

while \(f + g = (1, 1, 0, \ldots, 0)\) and so

\begin{equation*} \|f + g\|_{1/2} = \bigl( |1|^{1/2} + |1|^{1/2} \bigr)^2 = 2^2 = 4 . \end{equation*}

Thus \(\|f + g\|_{1/2} = 4 > 2 = \|f\|_{1/2} + \|g\|_{1/2}\), so the triangle inequality in 6.33 fails. Hence \(\|\cdot\|_{1/2}\) is not a norm on \(\mathbf{F}^n\) when \(n > 1\).

Problem (6C.4)

Prove that each Cauchy sequence in a normed vector space is bounded (meaning that there is a real number that is greater than the norm of every element in the Cauchy sequence).

Solution

Suppose \(V\) is a normed vector space and \(f_1, f_2, \ldots\) is a Cauchy sequence in \(V\). As recorded in Section 6C, in the context of a normed vector space the Cauchy condition (6.12) states that for every \(\varepsilon > 0\) there exists \(n \in \mathbf{Z}^+\) such that \(\|f_j - f_k\| < \varepsilon\) for all integers \(j \ge n\) and \(k \ge n\).

Apply this with \(\varepsilon = 1\): there exists \(n \in \mathbf{Z}^+\) such that

\begin{equation*} \|f_j - f_k\| < 1 \quad \text{for all integers } j \ge n \text{ and } k \ge n . \end{equation*}

In particular, if \(k \ge n\) then \(\|f_k - f_n\| < 1\), and hence the triangle inequality gives

\begin{equation*} \|f_k\| = \|(f_k - f_n) + f_n\| \le \|f_k - f_n\| + \|f_n\| < 1 + \|f_n\| . \end{equation*}

Thus all terms of the sequence from the \(n^{\text{th}}\) onwards have norm less than \(1 + \|f_n\|\).

Only finitely many terms remain, namely \(f_1, \ldots, f_{n-1}\) (there are none of these if \(n = 1\)). Define

\begin{equation*} M = \max \bigl\{ \|f_1\|, \ldots, \|f_{n-1}\|, \, 1 + \|f_n\| \bigr\}, \end{equation*}

which is a well-defined real number because it is the maximum of a finite nonempty set of real numbers.

We claim that the real number \(M + 1\) is greater than the norm of every element of the sequence. Indeed, if \(k < n\) then \(\|f_k\| \le M < M + 1\), and if \(k \ge n\) then \(\|f_k\| < 1 + \|f_n\| \le M < M + 1\). Hence

\begin{equation*} \|f_k\| < M + 1 \quad \text{for every } k \in \mathbf{Z}^+ , \end{equation*}

which is exactly the assertion that the Cauchy sequence \(f_1, f_2, \ldots\) is bounded.

This result is used in the proof of 6.47, where the boundedness of a Cauchy sequence \(T_1, T_2, \ldots\) in \(\mathcal{B}(V,W)\) guarantees that \(\sup\{\|T_k\| : k \in \mathbf{Z}^+\} < \infty\).

Problem (6C.5)

Show that if \(n \in \mathbf{Z}^+\), then \(\mathbf{F}^n\) is a Banach space with both the norms used in the first bullet point of Example 6.34.

Solution

The two norms in the first bullet point of 6.34 are defined on \(\mathbf{F}^n\) by

\begin{equation*} \|(a_1, \ldots, a_n)\|_1 = |a_1| + \cdots + |a_n| \qquad \text{and} \qquad \|(a_1, \ldots, a_n)\|_\infty = \max\{|a_1|, \ldots, |a_n|\} . \end{equation*}

Step 1: both are norms. Positive definiteness is clear in each case, since each of \(|a_1| + \cdots + |a_n|\) and \(\max\{|a_1|, \ldots, |a_n|\}\) equals \(0\) exactly when \(a_1 = \cdots = a_n = 0\). Homogeneity follows from \(|\alpha a_m| = |\alpha| |a_m|\), which gives \(\sum_m |\alpha a_m| = |\alpha| \sum_m |a_m|\) and \(\max_m |\alpha a_m| = |\alpha| \max_m |a_m|\). For the triangle inequality, if \(a = (a_1, \ldots, a_n)\) and \(b = (b_1, \ldots, b_n)\), then \(|a_m + b_m| \le |a_m| + |b_m|\) for each \(m\); summing over \(m\) gives \(\|a + b\|_1 \le \|a\|_1 + \|b\|_1\), and taking the maximum over \(m\) gives

\begin{equation*} \|a + b\|_\infty \le \max_m \bigl( |a_m| + |b_m| \bigr) \le \max_m |a_m| + \max_m |b_m| = \|a\|_\infty + \|b\|_\infty . \end{equation*}

Step 2: the two norms are comparable. For every \(a = (a_1, \ldots, a_n) \in \mathbf{F}^n\) and every \(m \in \{1, \ldots, n\}\),

\begin{equation*} |a_m| \le \|a\|_\infty \le \|a\|_1 \le n \|a\|_\infty , \end{equation*}

where the middle inequality holds because the maximum of the nonnegative numbers \(|a_1|, \ldots, |a_n|\) is one of the terms in the sum \(|a_1| + \cdots + |a_n|\), and the last inequality holds because each of the \(n\) terms of that sum is at most \(\|a\|_\infty\).

Step 3: completeness. Let \(\|\cdot\|\) denote either \(\|\cdot\|_1\) or \(\|\cdot\|_\infty\), and suppose \(f_1, f_2, \ldots\) is a Cauchy sequence in \(\mathbf{F}^n\) with respect to \(\|\cdot\|\). Write

\begin{equation*} f_k = (a_{k,1}, \ldots, a_{k,n}) \quad \text{for each } k \in \mathbf{Z}^+ . \end{equation*}

Fix \(m \in \{1, \ldots, n\}\). By Step 2, for all \(j, k \in \mathbf{Z}^+\) we have

\begin{equation*} |a_{j,m} - a_{k,m}| \le \|f_j - f_k\|_\infty \le \|f_j - f_k\| , \end{equation*}

the last inequality holding trivially if \(\|\cdot\| = \|\cdot\|_\infty\) and by Step 2 if \(\|\cdot\| = \|\cdot\|_1\). Because \(f_1, f_2, \ldots\) is Cauchy with respect to \(\|\cdot\|\), this shows that \(a_{1,m}, a_{2,m}, \ldots\) is a Cauchy sequence in \(\mathbf{F}\). Since \(\mathbf{F}\) (which is \(\mathbf{R}\) or \(\mathbf{C}\)) with the absolute value is a complete metric space, there exists \(b_m \in \mathbf{F}\) with

\begin{equation*} \lim_{k \to \infty} a_{k,m} = b_m . \end{equation*}

Let \(b = (b_1, \ldots, b_n) \in \mathbf{F}^n\). We claim that \(\lim_{k \to \infty} f_k = b\) with respect to \(\|\cdot\|_1\) (and hence, by Step 2, also with respect to \(\|\cdot\|_\infty\)). Suppose \(\varepsilon > 0\). For each \(m \in \{1, \ldots, n\}\) there exists \(N_m \in \mathbf{Z}^+\) such that

\begin{equation*} |a_{k,m} - b_m| < \frac{\varepsilon}{n} \quad \text{for all } k \ge N_m . \end{equation*}

Let \(N = \max\{N_1, \ldots, N_n\}\), which is a positive integer because it is the maximum of finitely many positive integers. If \(k \ge N\), then

\begin{equation*} \|f_k - b\|_1 = \sum_{m=1}^{n} |a_{k,m} - b_m| < n \cdot \frac{\varepsilon}{n} = \varepsilon , \end{equation*}

and also \(\|f_k - b\|_\infty \le \|f_k - b\|_1 < \varepsilon\) by Step 2. Hence \(\lim_{k \to \infty} \|f_k - b\| = 0\) for either choice of \(\|\cdot\|\), which means \(\lim_{k \to \infty} f_k = b\) in \((\mathbf{F}^n, \|\cdot\|)\).

Thus every Cauchy sequence in \((\mathbf{F}^n, \|\cdot\|)\) converges to an element of \(\mathbf{F}^n\), so \((\mathbf{F}^n, \|\cdot\|)\) is complete. Being a complete normed vector space, it is a Banach space (see 6.37), for both \(\|\cdot\| = \|\cdot\|_1\) and \(\|\cdot\| = \|\cdot\|_\infty\).

Problem (6C.6)

Suppose \(X\) is a nonempty set and \(b(X)\) is the vector space of bounded functions from \(X\) to \(\mathbf{F}\). Prove that if \(\|\cdot\|\) is defined on \(b(X)\) by \(\|f\| = \sup_X |f|\), then \(b(X)\) is a Banach space.

Solution

Here \(b(X)\) is the subspace of \(\mathbf{F}^X\) consisting of the bounded functions from \(X\) to \(\mathbf{F}\), and for \(f \in b(X)\) we set \(\|f\| = \sup\{|f(x)| : x \in X\}\), which is a finite nonnegative real number because \(f\) is bounded and \(X\) is nonempty (see the third bullet point of 6.34).

Step 1: \(\|\cdot\|\) is a norm on \(b(X)\). If \(\|f\| = 0\), then \(|f(x)| \le 0\) for every \(x \in X\), so \(f(x) = 0\) for every \(x \in X\), that is, \(f = 0\); conversely \(\|0\| = 0\). If \(\alpha \in \mathbf{F}\) and \(f \in b(X)\), then

\begin{equation*} \|\alpha f\| = \sup_{x \in X} |\alpha| \, |f(x)| = |\alpha| \sup_{x \in X} |f(x)| = |\alpha| \, \|f\| . \end{equation*}

If \(f, g \in b(X)\), then for every \(x \in X\) we have

\begin{equation*} |(f + g)(x)| \le |f(x)| + |g(x)| \le \|f\| + \|g\| , \end{equation*}

and taking the supremum over \(x \in X\) gives \(\|f + g\| \le \|f\| + \|g\|\). Thus all three conditions of 6.33 hold.

Step 2: construction of a candidate limit. Suppose \(f_1, f_2, \ldots\) is a Cauchy sequence in \(b(X)\). Fix \(x \in X\). For all \(j, k \in \mathbf{Z}^+\),

\begin{equation*} |f_j(x) - f_k(x)| \le \sup_{y \in X} |f_j(y) - f_k(y)| = \|f_j - f_k\| , \end{equation*}

so \(f_1(x), f_2(x), \ldots\) is a Cauchy sequence of elements of \(\mathbf{F}\). Because \(\mathbf{F}\) is complete, this sequence converges; define

\begin{equation*} f(x) = \lim_{k \to \infty} f_k(x) . \end{equation*}

This defines a function \(f \colon X \to \mathbf{F}\).

Step 3: \(f\) is bounded. By Exercise 4 in this section, the Cauchy sequence \(f_1, f_2, \ldots\) is bounded: there exists \(M \in \mathbf{R}\) with \(\|f_k\| < M\) for every \(k \in \mathbf{Z}^+\). Hence \(|f_k(x)| < M\) for all \(x \in X\) and all \(k \in \mathbf{Z}^+\). Because the absolute value function on \(\mathbf{F}\) is continuous, for each \(x \in X\) we have

\begin{equation*} |f(x)| = \lim_{k \to \infty} |f_k(x)| \le M . \end{equation*}

Thus \(f\) is bounded, so \(f \in b(X)\) and \(\|f\| \le M\).

Step 4: \(f_k \to f\) in the norm of \(b(X)\). Suppose \(\varepsilon > 0\). Because \(f_1, f_2, \ldots\) is Cauchy, there exists \(n \in \mathbf{Z}^+\) such that

\begin{equation*} \|f_j - f_k\| < \frac{\varepsilon}{2} \quad \text{for all integers } j \ge n \text{ and } k \ge n . \end{equation*}

Now fix \(k \ge n\) and \(x \in X\). For every \(j \ge n\) we have \(|f_j(x) - f_k(x)| \le \|f_j - f_k\| < \frac{\varepsilon}{2}\). Letting \(j \to \infty\) and using \(\lim_{j \to \infty} f_j(x) = f(x)\) together with the continuity of the absolute value, we get

\begin{equation*} |f(x) - f_k(x)| = \lim_{j \to \infty} |f_j(x) - f_k(x)| \le \frac{\varepsilon}{2} . \end{equation*}

This holds for every \(x \in X\), so taking the supremum over \(x \in X\) gives

\begin{equation*} \|f - f_k\| \le \frac{\varepsilon}{2} < \varepsilon \quad \text{for every } k \ge n . \end{equation*}

(The function \(f - f_k\) lies in \(b(X)\) because \(f\) and \(f_k\) both do.) Hence \(\lim_{k \to \infty} \|f_k - f\| = 0\), which means \(\lim_{k \to \infty} f_k = f\) in \(b(X)\).

Therefore every Cauchy sequence in \(b(X)\) converges to an element of \(b(X)\), so \(b(X)\) is complete. Being a complete normed vector space, \(b(X)\) is a Banach space (see 6.37).

Note that the norm \(\|\cdot\|\) on \(b(X)\) is the norm of uniform convergence: Step 4 shows precisely that a Cauchy sequence in \(b(X)\) converges pointwise to a bounded function and that the convergence is in fact uniform on \(X\).

Problem (6C.7)

Show that \(\ell^1\) with the norm defined by \(\|(a_1, a_2, \ldots)\|_\infty = \sup_{k \in \mathbf{Z}^+} |a_k|\) is not a Banach space.

Solution

Recall that \(\ell^1\) is the vector space of sequences \((a_1, a_2, \ldots)\) of elements of \(\mathbf{F}\) such that \(\sum_{k=1}^{\infty} |a_k| < \infty\) (see the last bullet point of 6.32). For a sequence \(h = (c_1, c_2, \ldots)\) we write \(h(m) = c_m\) for its \(m^{\text{th}}\) coordinate. What is to be proved is the assertion made in the second bullet point of Example 6.39.

First note that \(\|\cdot\|_\infty\) really is a norm on \(\ell^1\). It is well defined and finite there, because if \((a_1, a_2, \ldots) \in \ell^1\) then \(|a_m| \le \sum_{k=1}^{\infty} |a_k| < \infty\) for each \(m\), so the sequence is bounded. Positive definiteness holds because \(\sup_k |a_k| = 0\) forces every \(a_k\) to be \(0\); homogeneity holds because \(\sup_k |\alpha a_k| = |\alpha| \sup_k |a_k|\); and the triangle inequality holds because \(|a_k + b_k| \le |a_k| + |b_k| \le \|(a_1,a_2,\ldots)\|_\infty + \|(b_1,b_2,\ldots)\|_\infty\) for each \(k\), after which we take the supremum over \(k\). Thus \((\ell^1, \|\cdot\|_\infty)\) is a normed vector space, and to show it is not a Banach space we must exhibit a Cauchy sequence in \(\ell^1\) that does not converge (with respect to \(\|\cdot\|_\infty\)) to an element of \(\ell^1\).

For \(k \in \mathbf{Z}^+\) define

\begin{equation*} f_k = \Bigl( 1, \tfrac{1}{2}, \tfrac{1}{3}, \ldots, \tfrac{1}{k}, 0, 0, \ldots \Bigr), \end{equation*}

so the \(m^{\text{th}}\) coordinate of \(f_k\) is \(\frac{1}{m}\) if \(m \le k\) and is \(0\) if \(m > k\). Each \(f_k\) has only finitely many nonzero coordinates, so \(\sum_{m=1}^{\infty} |f_k(m)| = \sum_{m=1}^{k} \frac{1}{m} < \infty\) and hence \(f_k \in \ell^1\).

The sequence \(f_1, f_2, \ldots\) is Cauchy with respect to \(\|\cdot\|_\infty\). Indeed, if \(j > k \ge 1\), then the coordinates of \(f_j - f_k\) are \(0\) except in positions \(k+1, \ldots, j\), where they equal \(\frac{1}{k+1}, \ldots, \frac{1}{j}\); the largest of these in absolute value is \(\frac{1}{k+1}\), so

\begin{equation*} \|f_j - f_k\|_\infty = \frac{1}{k+1} . \end{equation*}

Hence for all \(j, k \ge n\) we have \(\|f_j - f_k\|_\infty \le \frac{1}{n+1}\) (this also covers \(j = k\), where the norm is \(0\), and \(j < k\) by symmetry). Given \(\varepsilon > 0\), choose \(n \in \mathbf{Z}^+\) with \(n > \frac{1}{\varepsilon}\); then \(\|f_j - f_k\|_\infty \le \frac{1}{n+1} < \varepsilon\) for all \(j \ge n\) and \(k \ge n\). So \(f_1, f_2, \ldots\) is a Cauchy sequence in \((\ell^1, \|\cdot\|_\infty)\).

Now suppose, seeking a contradiction, that this sequence converges in \((\ell^1, \|\cdot\|_\infty)\), say

\begin{equation*} \lim_{k \to \infty} \|f_k - g\|_\infty = 0 \quad \text{for some } g = (b_1, b_2, \ldots) \in \ell^1 . \end{equation*}

Fix \(m \in \mathbf{Z}^+\). For every \(k \ge m\) the \(m^{\text{th}}\) coordinate of \(f_k\) equals \(\frac{1}{m}\), so

\begin{equation*} \Bigl| \tfrac{1}{m} - b_m \Bigr| = |f_k(m) - b_m| \le \|f_k - g\|_\infty . \end{equation*}

Letting \(k \to \infty\), the right side tends to \(0\) while the left side does not depend on \(k\); hence \(\bigl| \frac{1}{m} - b_m \bigr| = 0\), that is, \(b_m = \frac{1}{m}\). Because \(m \in \mathbf{Z}^+\) was arbitrary,

\begin{equation*} g = \Bigl( 1, \tfrac{1}{2}, \tfrac{1}{3}, \ldots \Bigr) . \end{equation*}

But then

\begin{equation*} \sum_{m=1}^{\infty} |b_m| = \sum_{m=1}^{\infty} \frac{1}{m} = \infty \end{equation*}

because the harmonic series diverges, so \(g \notin \ell^1\), contradicting the assumption that \(g \in \ell^1\).

Therefore the Cauchy sequence \(f_1, f_2, \ldots\) has no limit in \(\ell^1\) with respect to \(\|\cdot\|_\infty\). Hence \((\ell^1, \|\cdot\|_\infty)\) is not complete, and so it is not a Banach space.

The failure is caused by using the wrong norm rather than by any defect of \(\ell^1\) itself: with its own norm \(\|(a_1, a_2, \ldots)\|_1 = \sum_{k=1}^{\infty} |a_k|\), the space \(\ell^1\) is a Banach space (see the second bullet point of 6.38 and Exercise 8 in this section). Indeed the sequence \(f_1, f_2, \ldots\) constructed above is not Cauchy with respect to \(\|\cdot\|_1\): if \(j > k\) then \(\|f_j - f_k\|_1 = \sum_{m=k+1}^{j} \frac{1}{m}\), and because the harmonic series diverges, for each \(k \in \mathbf{Z}^+\) there exists \(j > k\) with \(\sum_{m=k+1}^{j} \frac{1}{m} > 1\). Hence the Cauchy condition fails for \(\varepsilon = 1\).

Problem (6C.8)

Show that \(\ell^1\) with the norm defined by \(\|(a_1, a_2, \ldots)\|_1 = \sum_{k=1}^{\infty} |a_k|\) is a Banach space.

Solution

Recall that \(\ell^1\) is the vector space of sequences \((a_1, a_2, \ldots)\) of elements of \(\mathbf{F}\) such that \(\sum_{k=1}^{\infty} |a_k| < \infty\), with coordinatewise addition and scalar multiplication.

Step 1: \(\|\cdot\|_1\) is a norm. If \(\|(a_1,a_2,\ldots)\|_1 = 0\), then \(\sum_{k=1}^\infty |a_k| = 0\); since every term is nonnegative, \(|a_k| = 0\) for each \(k\), so \((a_1,a_2,\ldots) = 0\). Conversely \(\|0\|_1 = 0\). For \(\alpha \in \mathbf{F}\) we have \(\sum_{k=1}^\infty |\alpha a_k| = |\alpha| \sum_{k=1}^\infty |a_k|\), giving homogeneity. For the triangle inequality, note that for each \(n \in \mathbf{Z}^+\),

\begin{equation*} \sum_{k=1}^{n} |a_k + b_k| \le \sum_{k=1}^{n} |a_k| + \sum_{k=1}^{n} |b_k| \le \|(a_1,a_2,\ldots)\|_1 + \|(b_1,b_2,\ldots)\|_1 ; \end{equation*}

letting \(n \to \infty\) gives \(\|(a_1,a_2,\ldots) + (b_1,b_2,\ldots)\|_1 \le \|(a_1,a_2,\ldots)\|_1 + \|(b_1,b_2,\ldots)\|_1\). Thus \(\|\cdot\|_1\) is a norm on \(\ell^1\) (and, incidentally, this shows \(\ell^1\) is closed under addition).

Step 2: completeness. Suppose \(f^1, f^2, \ldots\) is a Cauchy sequence in \(\ell^1\), where

\begin{equation*} f^n = (a_1^n, a_2^n, \ldots). \end{equation*}

Fix \(k \in \mathbf{Z}^+\). For all \(m, n \in \mathbf{Z}^+\) we have

\begin{equation*} |a_k^n - a_k^m| \le \sum_{j=1}^{\infty} |a_j^n - a_j^m| = \|f^n - f^m\|_1 . \end{equation*}

Hence \(a_k^1, a_k^2, \ldots\) is a Cauchy sequence in \(\mathbf{F}\). Because \(\mathbf{F}\) (either \(\mathbb{R}\) or \(\mathbb{C}\)) is complete, this sequence converges; call its limit \(a_k\). Let \(f = (a_1, a_2, \ldots)\).

Now let \(\varepsilon > 0\). Because \(f^1, f^2, \ldots\) is Cauchy, there exists \(N \in \mathbf{Z}^+\) such that

\begin{equation*} \|f^n - f^m\|_1 < \varepsilon \quad \text{for all } m, n \ge N . \end{equation*}

Fix \(n \ge N\) and fix \(M \in \mathbf{Z}^+\). For every \(m \ge N\),

\begin{equation*} \sum_{k=1}^{M} |a_k^n - a_k^m| \le \|f^n - f^m\|_1 < \varepsilon . \end{equation*}

The left side is a finite sum, so we may let \(m \to \infty\) term by term, obtaining

\begin{equation*} \sum_{k=1}^{M} |a_k^n - a_k| \le \varepsilon . \end{equation*}

This bound holds for every \(M \in \mathbf{Z}^+\), and the partial sums are increasing in \(M\), so

\begin{equation*} \sum_{k=1}^{\infty} |a_k^n - a_k| \le \varepsilon \quad \text{for all } n \ge N . \end{equation*}

In particular \(f^N - f \in \ell^1\). Because \(\ell^1\) is a vector space and \(f^N \in \ell^1\), we conclude that

\begin{equation*} f = f^N - (f^N - f) \in \ell^1 . \end{equation*}

Moreover the displayed inequality says exactly that \(\|f^n - f\|_1 \le \varepsilon\) for all \(n \ge N\). Since \(\varepsilon > 0\) was arbitrary, \(\lim_{n \to \infty} \|f^n - f\|_1 = 0\), which means \(f^n \to f\) in \(\ell^1\).

Thus every Cauchy sequence in \(\ell^1\) converges to an element of \(\ell^1\), so \(\ell^1\) with \(\|\cdot\|_1\) is a Banach space (6.37).

Problem (6C.9)

Show that the vector space \(C([0,1])\) of continuous functions from \([0,1]\) to \(\mathbf{F}\) with the norm defined by \(\|f\| = \int_0^1 |f|\) is not a Banach space.

Solution

First note that the norm is well defined and finite: a continuous function on \([0,1]\) is Borel measurable and bounded, so \(\int_0^1 |f| < \infty\) for every \(f \in C([0,1])\). That \(\|\cdot\|\) satisfies homogeneity and the triangle inequality is immediate from the corresponding properties of the integral; the positive definite property follows from the Lemma below. We exhibit a Cauchy sequence in \(C([0,1])\) that does not converge in this norm.

First a small lemma that we will use twice.

Lemma. If \(g : [a,b] \to \mathbf{F}\) is continuous and \(\int_a^b |g| = 0\) (with \(a < b\)), then \(g = 0\) on \([a,b]\).

Proof of lemma. Suppose \(g(x_0) \ne 0\) for some \(x_0 \in [a,b]\). Let \(c = |g(x_0)| > 0\). By continuity of \(|g|\) there is an interval \(J \subseteq [a,b]\) of positive length \(\delta\) containing \(x_0\) on which \(|g| > c/2\). Since \(|g| \ge 0\) everywhere, \(\int_a^b |g| \ge \int_J |g| \ge c\delta/2 > 0\), a contradiction.

The sequence. For \(k \in \mathbf{Z}^+\) with \(k \ge 3\), define \(f_k : [0,1] \to \mathbf{F}\) by

\begin{equation*} f_k(x) = \begin{cases} 0 & \text{if } 0 \le x \le \tfrac12, \\[2pt] k\left(x - \tfrac12\right) & \text{if } \tfrac12 \le x \le \tfrac12 + \tfrac1k, \\[2pt] 1 & \text{if } \tfrac12 + \tfrac1k \le x \le 1. \end{cases} \end{equation*}

The three formulas agree at the two junction points, so each \(f_k\) is continuous, i.e. \(f_k \in C([0,1])\). Also \(0 \le f_k \le 1\) everywhere.

The sequence is Cauchy. Suppose \(j, k \ge 3\) and put \(m = \min\{j,k\}\). Both \(f_j\) and \(f_k\) vanish on \([0, \tfrac12]\) and both equal \(1\) on \([\tfrac12 + \tfrac1m, 1]\), so \(f_j - f_k\) vanishes outside the interval \([\tfrac12, \tfrac12 + \tfrac1m]\), which has length \(\tfrac1m\). On that interval \(|f_j - f_k| \le 1\). Hence

\begin{equation*} \|f_j - f_k\| = \int_0^1 |f_j - f_k| \le \frac{1}{m} = \frac{1}{\min\{j,k\}} . \end{equation*}

Given \(\varepsilon > 0\), choosing \(n > 1/\varepsilon\) gives \(\|f_j - f_k\| < \varepsilon\) for all \(j, k \ge n\). Thus \(f_3, f_4, \ldots\) is a Cauchy sequence.

The sequence does not converge in \(C([0,1])\). Suppose, for a contradiction, that there is \(f \in C([0,1])\) with \(\|f_k - f\| \to 0\).

On \([0,\tfrac12]\) we have \(f_k = 0\) for every \(k\), so

\begin{equation*} \int_0^{1/2} |f| = \int_0^{1/2} |f - f_k| \le \int_0^1 |f - f_k| = \|f_k - f\| \xrightarrow{\ k \to \infty\ } 0 . \end{equation*}

The left side does not depend on \(k\), so \(\int_0^{1/2} |f| = 0\), and the Lemma gives \(f = 0\) on \([0, \tfrac12]\).

Now fix \(\delta\) with \(0 < \delta < \tfrac12\). For every \(k > 1/\delta\) we have \(f_k = 1\) on \([\tfrac12 + \delta, 1]\), so

\begin{equation*} \int_{1/2 + \delta}^{1} |f - 1| = \int_{1/2+\delta}^{1} |f - f_k| \le \|f_k - f\| \xrightarrow{\ k \to \infty\ } 0 . \end{equation*}

Again the left side does not depend on \(k\), so \(\int_{1/2+\delta}^{1} |f - 1| = 0\), and the Lemma (applied to the continuous function \(f - 1\) on \([\tfrac12 + \delta, 1]\)) gives \(f = 1\) on \([\tfrac12 + \delta, 1]\). Since \(\delta \in (0, \tfrac12)\) was arbitrary, \(f = 1\) on \((\tfrac12, 1]\).

But then \(\lim_{x \downarrow 1/2} f(x) = 1\) while \(f(\tfrac12) = 0\), contradicting the continuity of \(f\) at \(\tfrac12\).

Hence no such \(f\) exists: the Cauchy sequence \(f_3, f_4, \ldots\) has no limit in \(C([0,1])\). Therefore \(C([0,1])\) with the norm \(\|f\| = \int_0^1 |f|\) is not complete, so it is not a Banach space.

Problem (6C.10)

Suppose \(U\) is a subspace of a normed vector space \(V\) such that some open ball of \(V\) is contained in \(U\). Prove that \(U = V\).

Solution

Suppose \(h \in V\) and \(r > 0\) are such that

\begin{equation*} B(h, r) = \{g \in V : \|g - h\| < r\} \subseteq U . \end{equation*}

Step 1: the ball centered at \(0\) of radius \(r\) is contained in \(U\). First note \(\|h - h\| = 0 < r\), so \(h \in B(h,r) \subseteq U\). Now suppose \(g \in V\) with \(\|g\| < r\). Then \(\|(h + g) - h\| = \|g\| < r\), so \(h + g \in B(h,r) \subseteq U\). Since \(U\) is a subspace and both \(h + g\) and \(h\) lie in \(U\), we get

\begin{equation*} g = (h + g) - h \in U . \end{equation*}

Thus \(B(0, r) \subseteq U\).

Step 2: \(U = V\). Certainly \(U \subseteq V\). For the reverse inclusion, let \(f \in V\). If \(f = 0\) then \(f \in U\) because \(U\) is a subspace. So suppose \(f \ne 0\), which means \(\|f\| > 0\) (by the positive definite property of the norm). Set

\begin{equation*} \alpha = \frac{r}{2\|f\|} > 0 . \end{equation*}

Then by homogeneity of the norm,

\begin{equation*} \|\alpha f\| = \alpha \|f\| = \frac{r}{2} < r , \end{equation*}

so \(\alpha f \in B(0,r) \subseteq U\) by Step 1. Since \(U\) is closed under scalar multiplication and \(\alpha \ne 0\),

\begin{equation*} f = \frac{1}{\alpha} (\alpha f) \in U . \end{equation*}

Hence \(V \subseteq U\), and therefore \(U = V\).

Problem (6C.11)

Prove that the only subsets of a normed vector space \(V\) that are both open and closed are \(\emptyset\) and \(V\).

Solution

In any metric space, \(\emptyset\) and the whole space are open (the defining condition is vacuous for \(\emptyset\), and every ball is contained in the whole space), and each is the complement of the other, so both are also closed. Thus we must show there are no others.

Suppose \(E \subseteq V\) is both open and closed and that \(E \ne \emptyset\) and \(E \ne V\). Pick \(f \in E\) and \(g \in V \setminus E\). Note \(f \ne g\), so \(\|g - f\| > 0\).

Define \(\gamma : [0,1] \to V\) by

\begin{equation*} \gamma(t) = (1 - t) f + t g = f + t (g - f) . \end{equation*}

For \(s, t \in [0,1]\), homogeneity of the norm gives

\begin{equation*} \|\gamma(s) - \gamma(t)\| = \|(s - t)(g - f)\| = |s - t| \, \|g - f\| , \end{equation*}

so \(\gamma\) is continuous (indeed Lipschitz). Note \(\gamma(0) = f \in E\) and \(\gamma(1) = g \notin E\).

Let

\begin{equation*} A = \{t \in [0,1] : \gamma(t) \in E\} . \end{equation*}

Then \(0 \in A\), so \(A \ne \emptyset\), and \(A\) is bounded above by \(1\). Let \(c = \sup A\), so \(c \in [0,1]\).

Claim: \(c \in A\). By definition of supremum there is a sequence \(t_1, t_2, \ldots\) in \(A\) with \(t_k \to c\). Then

\begin{equation*} \|\gamma(t_k) - \gamma( c)\| = |t_k - c| \, \|g - f\| \to 0 , \end{equation*}

so \(\gamma(t_k) \to \gamma( c)\) in \(V\). Each \(\gamma(t_k)\) lies in \(E\) and \(E\) is closed, so \(\gamma( c) \in E\) by 6.9(e). Thus \(c \in A\).

Since \(\gamma(1) = g \notin E\), we have \(1 \notin A\), so \(c < 1\).

Contradiction. Because \(E\) is open and \(\gamma( c) \in E\), there exists \(\varepsilon > 0\) with \(B(\gamma( c), \varepsilon) \subseteq E\). Choose

\begin{equation*} t \in \left(c, \ \min\left\{1, \ c + \frac{\varepsilon}{\|g - f\|}\right\}\right) , \end{equation*}

which is a nonempty interval because \(c < 1\) and \(\varepsilon / \|g - f\| > 0\). For such \(t\),

\begin{equation*} \|\gamma(t) - \gamma( c)\| = (t - c)\|g - f\| < \varepsilon , \end{equation*}

so \(\gamma(t) \in B(\gamma( c), \varepsilon) \subseteq E\), meaning \(t \in A\). But \(t > c = \sup A\), a contradiction.

Therefore no such \(E\) exists, and the only subsets of \(V\) that are both open and closed are \(\emptyset\) and \(V\).

Problem (6C.12)

Suppose \(V\) is a normed vector space. Prove that the closure of each subspace of \(V\) is a subspace of \(V\).

Solution

Let \(U\) be a subspace of \(V\) and let \(\overline{U}\) denote its closure. We must show \(\overline{U}\) contains \(0\) and is closed under addition and scalar multiplication.

By 6.9(a), \(\overline{U}\) is exactly the set of \(f \in V\) for which there exists a sequence \(f_1, f_2, \ldots\) in \(U\) with \(\lim_{k \to \infty} f_k = f\); we use this characterization throughout. (Recall that in a normed vector space \(\lim_{k\to\infty} f_k = f\) means \(\lim_{k\to\infty} \|f_k - f\| = 0\).)

Contains \(0\). Since \(U\) is a subspace, \(0 \in U \subseteq \overline{U}\). (The inclusion \(U \subseteq \overline{U}\) is immediate from 6.7, since \(g \in B(g,\varepsilon) \cap U\) for \(g \in U\).) In particular \(\overline{U}\) is nonempty.

Closed under addition. Suppose \(f, g \in \overline{U}\). By 6.9(a) choose sequences \(f_1, f_2, \ldots\) and \(g_1, g_2, \ldots\) in \(U\) with \(f_k \to f\) and \(g_k \to g\). Since \(U\) is a subspace, \(f_k + g_k \in U\) for every \(k\). By the triangle inequality,

\begin{equation*} \|(f_k + g_k) - (f + g)\| = \|(f_k - f) + (g_k - g)\| \le \|f_k - f\| + \|g_k - g\| , \end{equation*}

and the right side tends to \(0\) as \(k \to \infty\). Hence \(f_k + g_k \to f + g\), and \(f_k + g_k \in U\) for each \(k\), so \(f + g \in \overline{U}\) by 6.9(a).

Closed under scalar multiplication. Suppose \(f \in \overline{U}\) and \(\alpha \in \mathbf{F}\). Choose a sequence \(f_1, f_2, \ldots\) in \(U\) with \(f_k \to f\). Since \(U\) is a subspace, \(\alpha f_k \in U\) for every \(k\). By homogeneity of the norm,

\begin{equation*} \|\alpha f_k - \alpha f\| = \|\alpha(f_k - f)\| = |\alpha| \, \|f_k - f\| \to 0 , \end{equation*}

so \(\alpha f_k \to \alpha f\). Hence \(\alpha f \in \overline{U}\) by 6.9(a).

Therefore \(\overline{U}\) is a subspace of \(V\).

Problem (6C.13)

Suppose \(U\) is a normed vector space. Let \(d\) be the metric on \(U\) defined by \(d(f,g) = \|f - g\|\) for \(f, g \in U\). Let \(V\) be the complete metric space constructed in Exercise 16 in Section 6A.

(a) Show that the set \(V\) is a vector space under natural operations of addition and scalar multiplication.

(b) Show that there is a natural way to make \(V\) into a normed vector space and that with this norm, \(V\) is a Banach space.

(c) Explain why (b) shows that every normed vector space is a subspace of some Banach space.

Solution

Recall the construction of Exercise 16 in Section 6A. Let \(W\) be the set of all Cauchy sequences of elements of \(U\). For \((f_1, f_2, \ldots), (g_1, g_2, \ldots) \in W\) we write \((f_1, f_2, \ldots) \equiv (g_1, g_2, \ldots)\) to mean \(\lim_{k \to \infty} d(f_k, g_k) = 0\); this is an equivalence relation on \(W\). The set \(V\) is the set of equivalence classes, with \((f_1, f_2, \ldots)^{\wedge}\) denoting the class of \((f_1, f_2, \ldots)\), and

\begin{equation*} d_V\big((f_1, f_2, \ldots)^{\wedge}, (g_1, g_2, \ldots)^{\wedge}\big) = \lim_{k \to \infty} d(f_k, g_k) \end{equation*}

is a well-defined metric making \((V, d_V)\) a complete metric space. In our setting \(d(f,g) = \|f - g\|\).

(a) \(V\) is a vector space.

Define, for \((f_1, f_2, \ldots), (g_1, g_2, \ldots) \in W\) and \(\alpha \in \mathbf{F}\),

\begin{equation*} (f_1, f_2, \ldots)^{\wedge} + (g_1, g_2, \ldots)^{\wedge} = (f_1 + g_1, f_2 + g_2, \ldots)^{\wedge}, \qquad \alpha (f_1, f_2, \ldots)^{\wedge} = (\alpha f_1, \alpha f_2, \ldots)^{\wedge} . \end{equation*}

These land in \(W\). If \((f_k)\) and \((g_k)\) are Cauchy in \(U\), then for \(j, k \in \mathbf{Z}^+\),

\begin{equation*} \|(f_j + g_j) - (f_k + g_k)\| \le \|f_j - f_k\| + \|g_j - g_k\| , \end{equation*}

so given \(\varepsilon > 0\), taking \(n\) that works for both sequences with \(\varepsilon/2\) shows \((f_k + g_k)\) is Cauchy. Similarly \(\|\alpha f_j - \alpha f_k\| = |\alpha| \|f_j - f_k\|\) shows \((\alpha f_k)\) is Cauchy.

These are well defined on equivalence classes. Suppose \((f_k) \equiv (f_k^{\prime})\) and \((g_k) \equiv (g_k^{\prime})\). Then

\begin{equation*} \|(f_k + g_k) - (f_k^{\prime} + g_k^{\prime})\| \le \|f_k - f_k^{\prime}\| + \|g_k - g_k^{\prime}\| \to 0 , \end{equation*}

so \((f_k + g_k) \equiv (f_k^{\prime} + g_k^{\prime})\). Also \(\|\alpha f_k - \alpha f_k^{\prime}\| = |\alpha| \|f_k - f_k^{\prime}\| \to 0\), so \((\alpha f_k) \equiv (\alpha f_k^{\prime})\).

Vector space axioms. All the axioms are verified termwise, because each is an identity in \(U\) holding at every index. For instance commutativity: \((f_k + g_k)^{\wedge} = (g_k + f_k)^{\wedge}\) since \(f_k + g_k = g_k + f_k\) in \(U\) for each \(k\); associativity, distributivity, and \(1 \cdot v = v\) likewise. The additive identity is \((0,0,\ldots)^{\wedge}\) (the constant zero sequence is Cauchy), since \((f_k + 0)^{\wedge} = (f_k)^{\wedge}\); and the additive inverse of \((f_k)^{\wedge}\) is \((-f_k)^{\wedge} = (-1)(f_k)^{\wedge}\), since \((f_k + (-f_k))^{\wedge} = (0,0,\ldots)^{\wedge}\). Hence \(V\) is a vector space over \(\mathbf{F}\).

(b) \(V\) is a Banach space.

Define \(\|\cdot\|_V\) on \(V\) by

\begin{equation*} \big\|(f_1, f_2, \ldots)^{\wedge}\big\|_V = \lim_{k \to \infty} \|f_k\| . \end{equation*}

The limit exists. For \(j, k \in \mathbf{Z}^+\) the triangle inequality gives \(\big| \|f_j\| - \|f_k\| \big| \le \|f_j - f_k\|\). Since \((f_k)\) is Cauchy in \(U\), the sequence \(\|f_1\|, \|f_2\|, \ldots\) is a Cauchy sequence of real numbers, hence converges by completeness of \(\mathbb{R}\).

Well defined. If \((f_k) \equiv (g_k)\), then \(\big| \|f_k\| - \|g_k\| \big| \le \|f_k - g_k\| \to 0\), so \(\lim \|f_k\| = \lim \|g_k\|\).

Note that this is exactly \(\|(f_k)^{\wedge}\|_V = d_V\big((f_k)^{\wedge}, (0,0,\ldots)^{\wedge}\big)\), so the definition is the natural one. More importantly, for all \((f_k)^{\wedge}, (g_k)^{\wedge} \in V\),

\begin{equation*} \big\|(f_k)^{\wedge} - (g_k)^{\wedge}\big\|_V = \big\|(f_k - g_k)^{\wedge}\big\|_V = \lim_{k \to \infty} \|f_k - g_k\| = \lim_{k \to \infty} d(f_k, g_k) = d_V\big((f_k)^{\wedge}, (g_k)^{\wedge}\big) , \end{equation*}

so the metric induced by \(\|\cdot\|_V\) is precisely \(d_V\).

It is a norm. Positive definiteness: \(\|(f_k)^{\wedge}\|_V = 0\) means \(\lim_k \|f_k\| = 0\), i.e. \(\lim_k d(f_k, 0) = 0\), i.e. \((f_k) \equiv (0,0,\ldots)\), i.e. \((f_k)^{\wedge}\) is the zero vector of \(V\); and conversely. Homogeneity: \(\|\alpha (f_k)^{\wedge}\|_V = \lim_k \|\alpha f_k\| = |\alpha| \lim_k \|f_k\| = |\alpha| \, \|(f_k)^{\wedge}\|_V\). Triangle inequality: since \(\|f_k + g_k\| \le \|f_k\| + \|g_k\|\) for each \(k\), taking limits (all three limits exist) gives

\begin{equation*} \big\|(f_k)^{\wedge} + (g_k)^{\wedge}\big\|_V \le \big\|(f_k)^{\wedge}\big\|_V + \big\|(g_k)^{\wedge}\big\|_V . \end{equation*}

Thus \((V, \|\cdot\|_V)\) is a normed vector space.

It is complete. By 6.36 the metric associated with the normed vector space \((V, \|\cdot\|_V)\) is \(d_V\), as computed above. By Exercise 16(c) in Section 6A, \((V, d_V)\) is a complete metric space. Hence \((V, \|\cdot\|_V)\) is a complete normed vector space, i.e. a Banach space (6.37).

(c) Define \(\varphi : U \to V\) by \(\varphi(f) = (f, f, f, \ldots)^{\wedge}\) (a constant sequence is trivially Cauchy). Then \(\varphi\) is linear: for \(f, g \in U\) and \(\alpha \in \mathbf{F}\),

\begin{equation*} \varphi(f + g) = (f + g, f + g, \ldots)^{\wedge} = (f,f,\ldots)^{\wedge} + (g,g,\ldots)^{\wedge} = \varphi(f) + \varphi(g), \end{equation*}

and \(\varphi(\alpha f) = (\alpha f, \alpha f, \ldots)^{\wedge} = \alpha \varphi(f)\), by the definitions of the operations in (a). Moreover \(\varphi\) preserves norms:

\begin{equation*} \|\varphi(f)\|_V = \lim_{k \to \infty} \|f\| = \|f\| . \end{equation*}

In particular \(\varphi(f) = 0\) forces \(\|f\| = 0\), so \(f = 0\); hence \(\varphi\) is injective. Being the image of a linear map, \(\varphi(U)\) is a subspace of \(V\), and \(\varphi\) is a norm-preserving linear bijection from \(U\) onto \(\varphi(U)\).

Thus \(U\) is isometrically isomorphic, as a normed vector space, to the subspace \(\varphi(U)\) of the Banach space \(V\). Identifying \(U\) with \(\varphi(U)\) (which changes nothing about its vector space structure or its norm), we conclude that every normed vector space is a subspace of some Banach space.

Problem (6C.14)

Suppose \(U\) is a subspace of a normed vector space \(V\). Suppose also that \(W\) is a Banach space and \(S : U \to W\) is a bounded linear map.

(a) Prove that there exists a unique continuous function \(T : \overline{U} \to W\) such that \(T|_U = S\).

(b) Prove that the function \(T\) in (a) is a bounded linear map from \(\overline{U}\) to \(W\) and \(\|T\| = \|S\|\).

(c) Give an example to show that (a) can fail if the assumption that \(W\) is a Banach space is replaced by the assumption that \(W\) is a normed vector space.

Solution

Throughout, \(\overline{U}\) is the closure of \(U\) in \(V\); by 6C.12 it is a subspace of \(V\), so the statements about linearity make sense.

Preliminary: \(\|Sh\| \le \|S\| \, \|h\|\) for all \(h \in U\). If \(h = 0\) then \(Sh = 0\) and both sides are \(0\). If \(h \ne 0\), then \(\|h/\|h\|\| = 1\), so by the definition 6.43 of \(\|S\|\) as a supremum over the closed unit ball of \(U\),

\begin{equation*} \frac{\|Sh\|}{\|h\|} = \left\| S\!\left( \frac{h}{\|h\|} \right) \right\| \le \|S\| , \end{equation*}

which gives the claim. Note \(\|S\| < \infty\) since \(S\) is bounded.

(a) Existence and uniqueness.

Uniqueness. Suppose \(T_1, T_2 : \overline{U} \to W\) are continuous with \(T_1|_U = T_2|_U = S\). Let \(f \in \overline{U}\). By 6.9(a) there is a sequence \(f_1, f_2, \ldots\) in \(U\) with \(f_k \to f\). By continuity and 6.11 (continuity is equivalent to preservation of limits of sequences),

\begin{equation*} T_1 f = \lim_{k \to \infty} T_1 f_k = \lim_{k \to \infty} S f_k = \lim_{k \to \infty} T_2 f_k = T_2 f . \end{equation*}

Hence \(T_1 = T_2\).

Existence. Let \(f \in \overline{U}\) and choose (by 6.9(a)) a sequence \(f_1, f_2, \ldots\) in \(U\) with \(f_k \to f\). Since this sequence converges, it is Cauchy (6.13). By linearity of \(S\) and the preliminary estimate,

\begin{equation*} \|S f_j - S f_k\| = \|S(f_j - f_k)\| \le \|S\| \, \|f_j - f_k\| , \end{equation*}

so \(S f_1, S f_2, \ldots\) is a Cauchy sequence in \(W\). Because \(W\) is a Banach space, this sequence converges; define

\begin{equation*} T f = \lim_{k \to \infty} S f_k . \end{equation*}

This is well defined. Suppose \(g_1, g_2, \ldots\) is another sequence in \(U\) with \(g_k \to f\). Then

\begin{equation*} \|S f_k - S g_k\| \le \|S\| \, \|f_k - g_k\| \le \|S\| \big( \|f_k - f\| + \|f - g_k\| \big) \to 0 , \end{equation*}

so \(\lim_k S f_k = \lim_k S g_k\). Thus \(Tf\) does not depend on the choice of sequence.

\(T|_U = S\). If \(f \in U\), take the constant sequence \(f_k = f\); then \(Tf = \lim_k Sf = Sf\).

\(T\) is continuous. Let \(f, g \in \overline{U}\) and choose sequences \(f_k \to f\) and \(g_k \to g\) in \(U\). Then \(f_k - g_k \to f - g\), and for each \(k\),

\begin{equation*} \|S f_k - S g_k\| \le \|S\| \, \|f_k - g_k\| . \end{equation*}

Now \(Sf_k \to Tf\) and \(Sg_k \to Tg\), so \(Sf_k - Sg_k \to Tf - Tg\) (since \(\|(Sf_k - Sg_k) - (Tf - Tg)\| \le \|Sf_k - Tf\| + \|Sg_k - Tg\| \to 0\)). By 6C.1 the norm is a continuous function, so letting \(k \to \infty\) on both sides yields

\begin{equation*} \|Tf - Tg\| \le \|S\| \, \|f - g\| . \end{equation*}

Thus \(T\) is Lipschitz with constant \(\|S\|\), hence continuous. This proves (a).

(b) \(T\) is a bounded linear map with \(\|T\| = \|S\|\).

Linearity. Let \(f, g \in \overline{U}\) and \(\alpha \in \mathbf{F}\). Choose sequences \(f_k \to f\) and \(g_k \to g\) in \(U\). Since \(U\) is a subspace, \(f_k + g_k \in U\) and \(\alpha f_k \in U\); moreover \(f_k + g_k \to f + g\) and \(\alpha f_k \to \alpha f\) (as in the proof of 6C.12). Therefore, using the definition of \(T\) along these sequences and the linearity of \(S\),

\begin{equation*} T(f + g) = \lim_{k \to \infty} S(f_k + g_k) = \lim_{k \to \infty} (S f_k + S g_k) = Tf + Tg , \end{equation*}

where the last equality holds because \(\|(Sf_k + Sg_k) - (Tf + Tg)\| \le \|Sf_k - Tf\| + \|Sg_k - Tg\| \to 0\). Similarly

\begin{equation*} T(\alpha f) = \lim_{k \to \infty} S(\alpha f_k) = \lim_{k \to \infty} \alpha S f_k = \alpha T f , \end{equation*}

because \(\|\alpha S f_k - \alpha T f\| = |\alpha| \, \|S f_k - T f\| \to 0\). Hence \(T\) is linear.

Boundedness and the norm. Taking \(g = 0\) in the Lipschitz estimate from (a) (note \(0 \in U \subseteq \overline{U}\) and \(T0 = S0 = 0\)) gives

\begin{equation*} \|Tf\| \le \|S\| \, \|f\| \quad \text{for all } f \in \overline{U} . \end{equation*}

Hence for \(f \in \overline{U}\) with \(\|f\| \le 1\) we get \(\|Tf\| \le \|S\|\), so \(\|T\| \le \|S\| < \infty\); in particular \(T\) is bounded.

Conversely, since \(U \subseteq \overline{U}\) and \(T|_U = S\),

\begin{equation*} \|S\| = \sup\{\|S f\| : f \in U, \ \|f\| \le 1\} = \sup\{\|T f\| : f \in U, \ \|f\| \le 1\} \le \sup\{\|T f\| : f \in \overline{U}, \ \|f\| \le 1\} = \|T\| , \end{equation*}

the inequality holding because the supremum is over a larger set. Combining the two inequalities gives \(\|T\| = \|S\|\).

(c) A counterexample when \(W\) is merely a normed vector space.

Let \(V = \ell^1\) with the norm \(\|\cdot\|_1\), which is a Banach space by 6C.8. Let

\begin{equation*} U = \{(a_1, a_2, \ldots) \in \ell^1 : a_k = 0 \text{ for all but finitely many } k \in \mathbf{Z}^+\} , \end{equation*}

which is clearly a subspace of \(\ell^1\).

\(U\) is dense in \(\ell^1\) and \(U \ne \ell^1\). Given \(a = (a_1, a_2, \ldots) \in \ell^1\), let \(a^{(n)} = (a_1, \ldots, a_n, 0, 0, \ldots) \in U\). Then

\begin{equation*} \|a - a^{(n)}\|_1 = \sum_{k = n+1}^{\infty} |a_k| \to 0 \quad \text{as } n \to \infty , \end{equation*}

because the tails of a convergent series of nonnegative terms tend to \(0\). Hence \(a \in \overline{U}\) by 6.9(a), so \(\overline{U} = \ell^1\). On the other hand \((1/2, 1/4, 1/8, \ldots) \in \ell^1 \setminus U\), so \(U \ne \ell^1\).

Let \(W = U\), equipped with the norm \(\|\cdot\|_1\) inherited from \(\ell^1\). Then \(W\) is a normed vector space. (It is not a Banach space: a complete subspace of a metric space is closed, since a sequence in it converging in \(\ell^1\) is Cauchy, hence converges in \(W\) to a limit that must equal the \(\ell^1\) limit by uniqueness of limits; but \(\overline{U} = \ell^1 \ne U\), so \(U\) is not closed.)

Let \(S : U \to W\) be the identity map \(Sa = a\). It is linear, and \(\|Sa\|_1 = \|a\|_1\), so \(\|S\| = 1 < \infty\); thus \(S\) is a bounded linear map.

No continuous extension exists. Suppose \(T : \overline{U} = \ell^1 \to W\) were continuous with \(T|_U = S\). Let \(f = (1/2, 1/4, 1/8, \ldots) \in \ell^1 \setminus U\), and let \(f_k = f^{(k)} \in U\) be its truncations, so \(f_k \to f\) in \(\ell^1\) as shown above. By continuity of \(T\) and 6.11,

\begin{equation*} T f = \lim_{k \to \infty} T f_k = \lim_{k \to \infty} S f_k = \lim_{k \to \infty} f_k , \end{equation*}

where the limit is taken in \(W\), i.e. with respect to \(\|\cdot\|_1\). Since the norm on \(W\) is the restriction of the norm on \(\ell^1\), convergence in \(W\) implies convergence in \(\ell^1\) to the same element. But \(f_k \to f\) in \(\ell^1\) and limits in a metric space are unique, so \(Tf = f\). This is a contradiction, because \(Tf \in W = U\) while \(f \notin U\).

Hence part (a) fails for this \(U\), \(V\), \(W\), and \(S\), showing that the hypothesis that \(W\) is a Banach space cannot be weakened to the hypothesis that \(W\) is a normed vector space.

Problem (6C.15)

For readers familiar with the quotient of a vector space and a subspace: Suppose \(V\) is a normed vector space and \(U\) is a subspace of \(V\). Define \(\|\cdot\|\) on \(V/U\) by

\begin{equation*} \|f + U\| = \inf\{\|f + g\| : g \in U\}. \end{equation*}

(a) Prove that \(\|\cdot\|\) is a norm on \(V/U\) if and only if \(U\) is a closed subspace of \(V\).

(b) Prove that if \(V\) is a Banach space and \(U\) is a closed subspace of \(V\), then \(V/U\) (with the norm defined above) is a Banach space.

(c) Prove that if \(U\) is a Banach space (with the norm it inherits from \(V\)) and \(V/U\) is a Banach space (with the norm defined above), then \(V\) is a Banach space.

Solution

Throughout, \(V/U\) denotes the quotient vector space, whose elements are the cosets \(f + U = \{f + g : g \in U\}\) for \(f \in V\), with the operations

\begin{equation*} (f + U) + (h + U) = (f + h) + U, \qquad \alpha(f + U) = (\alpha f) + U . \end{equation*}

Its zero element is \(0 + U = U\), and \(f + U = 0 + U\) if and only if \(f \in U\).

Preliminary observations: the formula is well defined, finite, and equals the distance to \(U\).

If \(f + U = f^{\prime} + U\), then \(\{f + g : g \in U\} = \{f^{\prime} + g : g \in U\}\) as subsets of \(V\) (both equal the coset), so the two infima defining \(\|f + U\|\) and \(\|f^{\prime} + U\|\) are infima of the same set of numbers. Hence \(\|\cdot\|\) is a well-defined function of the coset. Taking \(g = 0 \in U\) shows \(0 \le \|f + U\| \le \|f\| < \infty\), so \(\|\cdot\| : V/U \to [0, \infty)\).

Because \(g \mapsto -g\) is a bijection of \(U\) onto \(U\),

\begin{equation*} \|f + U\| = \inf\{\|f + g\| : g \in U\} = \inf\{\|f - g\| : g \in U\} = \operatorname{dist}(f, U), \end{equation*}

the distance from \(f\) to the set \(U\) in the metric of 6.36.

Part (a). By 6.33, a norm on \(V/U\) is a function from \(V/U\) to \([0, \infty)\) that is positive definite, homogeneous, and satisfies the triangle inequality. The preliminary remarks already show that \(\|\cdot\|\) maps \(V/U\) into \([0, \infty)\). We now check that homogeneity and the triangle inequality hold with no hypothesis on \(U\) beyond its being a subspace, so that the closedness of \(U\) is exactly what the positive definite condition needs.

Homogeneity. If \(\alpha = 0\), then \(\|\alpha(f + U)\| = \|0 + U\| = \inf\{\|g\| : g \in U\} = 0\) (take \(g = 0\)), which equals \(|\alpha| \, \|f + U\|\). If \(\alpha \ne 0\), then \(g \mapsto \alpha g\) is a bijection of \(U\) onto \(U\), so

\begin{equation*} \|\alpha(f + U)\| = \inf_{g \in U} \|\alpha f + g\| = \inf_{h \in U} \|\alpha f + \alpha h\| = \inf_{h \in U} |\alpha| \, \|f + h\| = |\alpha| \, \|f + U\| . \end{equation*}

Triangle inequality. Suppose \(f, h \in V\) and \(g_1, g_2 \in U\). Then \(g_1 + g_2 \in U\), so

\begin{equation*} \|(f + h) + U\| \le \|(f + h) + (g_1 + g_2)\| \le \|f + g_1\| + \|h + g_2\| . \end{equation*}

Taking the infimum over \(g_1 \in U\) and then over \(g_2 \in U\) gives \(\|(f + h) + U\| \le \|f + U\| + \|h + U\|\), which is the triangle inequality for the sum \((f + U) + (h + U)\).

Positive definiteness. Always \(\|0 + U\| = 0\), as noted above. So \(\|\cdot\|\) is a norm on \(V/U\) if and only if the implication

\begin{equation*} \|f + U\| = 0 \implies f + U = 0 + U \end{equation*}

holds for every \(f \in V\); that is, if and only if

\begin{equation*} \operatorname{dist}(f, U) = 0 \implies f \in U \quad \text{for every } f \in V. \end{equation*}

Now \(\operatorname{dist}(f, U) = 0\) if and only if for each \(k \in \mathbb{Z}^+\) there is \(g_k \in U\) with \(\|f - g_k\| < 1/k\), which happens if and only if \(f\) is the limit of a sequence in \(U\), which happens if and only if \(f \in \overline{U}\). Hence the displayed implication holds for all \(f \in V\) if and only if \(\overline{U} \subseteq U\), i.e. if and only if \(U\) is closed (the inclusion \(U \subseteq \overline{U}\) is automatic). This proves (a).

Part (b). Suppose \(V\) is a Banach space and \(U\) is a closed subspace. By (a), \(V/U\) is a normed vector space. We use 6.41: it suffices to show that \(\sum_{k=1}^{\infty}(f_k + U)\) converges in \(V/U\) whenever \(\sum_{k=1}^{\infty} \|f_k + U\| < \infty\).

So suppose \(f_1, f_2, \ldots \in V\) and \(\sum_{k=1}^{\infty} \|f_k + U\| < \infty\). For each \(k \in \mathbb{Z}^+\), the definition of the infimum provides \(g_k \in U\) with

\begin{equation*} \|f_k + g_k\| \le \|f_k + U\| + 2^{-k} . \end{equation*}

Then \(\sum_{k=1}^{\infty} \|f_k + g_k\| \le \sum_{k=1}^{\infty} \|f_k + U\| + 1 < \infty\). Because \(V\) is a Banach space, 6.41 implies that \(\sum_{k=1}^{\infty} (f_k + g_k)\) converges in \(V\); call its sum \(f\), so that

\begin{equation*} \Bigl\| \sum_{k=1}^{n} (f_k + g_k) - f \Bigr\| \longrightarrow 0 \quad \text{as } n \to \infty . \end{equation*}

For each \(n\) we have \(\sum_{k=1}^{n} g_k \in U\), so

\begin{equation*} \sum_{k=1}^{n} (f_k + U) = \Bigl(\sum_{k=1}^{n} f_k\Bigr) + U = \Bigl(\sum_{k=1}^{n} (f_k + g_k)\Bigr) + U , \end{equation*}

and therefore

\begin{equation*} \Bigl\| \sum_{k=1}^{n} (f_k + U) - (f + U) \Bigr\| = \Bigl\| \Bigl(\sum_{k=1}^{n} (f_k + g_k) - f\Bigr) + U \Bigr\| \le \Bigl\| \sum_{k=1}^{n} (f_k + g_k) - f \Bigr\| , \end{equation*}

where the inequality is the estimate \(\|h + U\| \le \|h\|\) from the preliminary remarks. The right side tends to \(0\), so \(\sum_{k=1}^{\infty}(f_k + U)\) converges in \(V/U\) (to \(f + U\)). By 6.41 again, \(V/U\) is a Banach space.

Part (c). Suppose \(U\) is complete in the norm it inherits from \(V\), and \(V/U\) is a Banach space. First note that \(U\) is closed in \(V\): if \(f \in \overline{U}\), pick \(g_k \in U\) with \(g_k \to f\) in \(V\); then \(g_1, g_2, \ldots\) is a Cauchy sequence in \(U\) (by 6.13), so completeness of \(U\) gives \(g \in U\) with \(g_k \to g\); limits in a metric space are unique, so \(f = g \in U\). Thus \(\overline{U} = U\) and, by (a), \(\|\cdot\|\) really is a norm on \(V/U\).

Now let \(f_1, f_2, \ldots\) be a Cauchy sequence in \(V\); we must produce a limit in \(V\). Since

\begin{equation*} \|(f_j + U) - (f_k + U)\| = \|(f_j - f_k) + U\| \le \|f_j - f_k\| , \end{equation*}

the sequence \(f_1 + U, f_2 + U, \ldots\) is Cauchy in \(V/U\). Because \(V/U\) is a Banach space, this sequence converges; write its limit as \(f + U\) for some \(f \in V\). Thus

\begin{equation*} \|(f_k - f) + U\| \longrightarrow 0 \quad \text{as } k \to \infty . \end{equation*}

For each \(k \in \mathbb{Z}^+\), the definition of the infimum provides \(h_k \in U\) with

\begin{equation*} \|f_k - f + h_k\| \le \|(f_k - f) + U\| + \tfrac{1}{k} . \end{equation*}

Set \(u_k = f_k + h_k\). The last display says \(\|u_k - f\| \le \|(f_k - f) + U\| + 1/k \to 0\), so \(u_1, u_2, \ldots\) converges to \(f\) in \(V\).

Consider now \(h_k = u_k - f_k\). For \(j, k \in \mathbb{Z}^+\),

\begin{equation*} \|h_j - h_k\| \le \|u_j - u_k\| + \|f_j - f_k\| . \end{equation*}

The sequence \(u_1, u_2, \ldots\) converges and hence is Cauchy (6.13), and \(f_1, f_2, \ldots\) is Cauchy by hypothesis; thus given \(\varepsilon > 0\) there is \(n\) with \(\|u_j - u_k\| < \varepsilon/2\) and \(\|f_j - f_k\| < \varepsilon/2\) for all \(j, k \ge n\), whence \(\|h_j - h_k\| < \varepsilon\). So \(h_1, h_2, \ldots\) is a Cauchy sequence in \(U\). Because \(U\) is complete, there exists \(h \in U\) with \(h_k \to h\).

Therefore

\begin{equation*} \|f_k - (f - h)\| = \|(u_k - h_k) - (f - h)\| \le \|u_k - f\| + \|h_k - h\| \longrightarrow 0 , \end{equation*}

so \(f_1, f_2, \ldots\) converges to \(f - h \in V\). Every Cauchy sequence in \(V\) converges, so \(V\) is a Banach space.

Problem (6C.16)

Suppose \(V\) and \(W\) are normed vector spaces with \(V \ne \{0\}\) and \(T : V \to W\) is a linear map.

(a) Show that \(\|T\| = \sup\{\|Tf\| : f \in V \text{ and } \|f\| < 1\}\).

(b) Show that \(\|T\| = \sup\{\|Tf\| : f \in V \text{ and } \|f\| = 1\}\).

(c) Show that \(\|T\| = \inf\{c \in [0, \infty) : \|Tf\| \le c\|f\| \text{ for all } f \in V\}\).

(d) Show that

\begin{equation*} \|T\| = \sup\Bigl\{ \frac{\|Tf\|}{\|f\|} : f \in V \text{ and } f \ne 0 \Bigr\}. \end{equation*}

Solution

Nothing here assumes \(T\) is bounded: all five quantities below are elements of \([0, \infty]\), and the asserted equalities are equalities in \([0, \infty]\). As usual, \(\sup\) of a set of nonnegative numbers that is unbounded above is \(\infty\), and \(\inf \emptyset = \infty\).

Write

\begin{equation*} \begin{aligned} A &= \sup\{\|Tf\| : f \in V,\ \|f\| \le 1\} \quad (\text{this is } \|T\|, \text{ by } 6.43), \\ B &= \sup\{\|Tf\| : f \in V,\ \|f\| < 1\}, \\ C &= \sup\{\|Tf\| : f \in V,\ \|f\| = 1\}, \\ D &= \inf\{c \in [0, \infty) : \|Tf\| \le c\|f\| \text{ for all } f \in V\}, \\ E &= \sup\{\|Tf\| / \|f\| : f \in V,\ f \ne 0\}. \end{aligned} \end{equation*}

The sets defining \(A\) and \(B\) contain \(0\), so they are nonempty. Because \(V \ne \{0\}\) there is some \(f \in V\) with \(f \ne 0\), and then \(f/\|f\|\) has norm \(1\) by homogeneity; so the sets defining \(C\) and \(E\) are nonempty too. Hence each of \(A, B, C, E\) is the supremum of a nonempty set of nonnegative real numbers and therefore lies in \([0, \infty]\). We must show \(A = B\), \(A = C\), \(A = D\), \(A = E\).

Two facts are used repeatedly. First, \(T0 = 0\) (linearity), so \(\|T0\| = 0\). Second, for \(f \in V\) and \(\alpha \in \mathbf{F}\), homogeneity of the norm on \(W\) and linearity of \(T\) give \(\|T(\alpha f)\| = \|\alpha \, Tf\| = |\alpha| \, \|Tf\|\).

Part (a): \(A = B\). Since \(\{f : \|f\| < 1\} \subseteq \{f : \|f\| \le 1\}\), we get \(B \le A\).

For the reverse inequality, suppose \(f \in V\) with \(\|f\| \le 1\), and let \(t \in (0,1)\). Then \(\|tf\| = t \|f\| \le t < 1\), so \(tf\) belongs to the set over which \(B\) is a supremum; hence

\begin{equation*} t \|Tf\| = \|T(tf)\| \le B . \end{equation*}

If \(B = \infty\) then certainly \(\|Tf\| \le B\). If \(B < \infty\), then \(\|Tf\| \le B/t\) for every \(t \in (0,1)\), and letting \(t \uparrow 1\) gives \(\|Tf\| \le B\). Either way \(\|Tf\| \le B\) for every \(f\) with \(\|f\| \le 1\), so \(A \le B\). Hence \(A = B\).

Part (b): \(A = C\). Since \(\{f : \|f\| = 1\} \subseteq \{f : \|f\| \le 1\}\), we get \(C \le A\).

For the reverse inequality, suppose \(f \in V\) with \(\|f\| \le 1\). If \(f = 0\) then \(\|Tf\| = 0 \le C\) (recall \(C \ge 0\)). If \(f \ne 0\), put \(u = f/\|f\|\), so \(\|u\| = 1\) and

\begin{equation*} \|Tf\| = \bigl\| T(\|f\| u) \bigr\| = \|f\| \, \|Tu\| \le \|Tu\| \le C , \end{equation*}

where the first inequality holds because \(\|f\| \le 1\) and \(\|Tu\|\) is a nonnegative real number. Thus \(\|Tf\| \le C\) for all \(f\) with \(\|f\| \le 1\), giving \(A \le C\). Hence \(A = C\).

Part (c): \(A = D\). Let

\begin{equation*} S = \{c \in [0, \infty) : \|Tf\| \le c \|f\| \text{ for all } f \in V\}, \end{equation*}

so that \(D = \inf S\).

First we show \(S \subseteq [A, \infty)\). Suppose \(c \in S\). If \(f \in V\) with \(\|f\| \le 1\), then \(\|Tf\| \le c\|f\| \le c\). Taking the supremum over such \(f\) gives \(A \le c\).

Next we show that \(\|Tf\| \le A \|f\|\) for every \(f \in V\). This is clear for \(f = 0\), since both sides are \(0\). If \(f \ne 0\), then \(u = f/\|f\|\) satisfies \(\|u\| = 1 \le 1\), so \(\|Tu\| \le A\), and hence

\begin{equation*} \|Tf\| = \|f\| \, \|Tu\| \le A \|f\| . \end{equation*}

Now conclude. If \(A < \infty\), the previous paragraph exhibits \(A\) as one of the numbers \(c \in [0, \infty)\) satisfying \(\|Tf\| \le c\|f\|\) for all \(f \in V\), so \(A \in S\); the paragraph before it shows every element of \(S\) is at least \(A\). Hence \(\inf S = A\), i.e. \(D = A\). If \(A = \infty\), then no \(c \in [0, \infty)\) satisfies \(c \ge A\), so the first step forces \(S = \emptyset\) and thus \(D = \inf \emptyset = \infty = A\). In both cases \(D = A\).

(The middle step just proved is worth recording on its own: \(\|Tf\| \le \|T\| \, \|f\|\) for all \(f \in V\), and by the first step \(\|T\|\) is the smallest constant \(c\) for which \(\|Tf\| \le c\|f\|\) holds for all \(f \in V\).)

Part (d): \(A = E\). We show that the two sets

\begin{equation*} P = \{\|Tf\|/\|f\| : f \in V,\ f \ne 0\} \quad \text{and} \quad Q = \{\|Tu\| : u \in V,\ \|u\| = 1\} \end{equation*}

are equal, so that \(E = \sup P = \sup Q = C\), and then \(E = C = A\) by part (b).

If \(u \in V\) with \(\|u\| = 1\), then \(u \ne 0\) (positive definiteness) and \(\|Tu\| = \|Tu\|/\|u\|\), so \(Q \subseteq P\). Conversely, if \(f \in V\) with \(f \ne 0\), put \(u = f/\|f\|\); then \(\|u\| = \|f\|/\|f\| = 1\) and

\begin{equation*} \|Tu\| = \frac{1}{\|f\|} \|Tf\| = \frac{\|Tf\|}{\|f\|}, \end{equation*}

so \(P \subseteq Q\). Hence \(P = Q\) and \(E = C = A\), which is the assertion of (d).

Finally, a remark on the hypothesis \(V \ne \{0\}\): it is needed only for (b) and (d). If \(V = \{0\}\), then \(\|T\| = \sup\{\|T0\|\} = 0\), while the sets \(\{f \in V : \|f\| = 1\}\) and \(\{f \in V : f \ne 0\}\) are both empty, so (b) and (d) assert that \(0\) equals the supremum of the empty set. Parts (a) and (c) remain true when \(V = \{0\}\): the set in (a) is \(\{0\}\), whose image under \(T\) has supremum \(0\), and the set in (c) is all of \([0, \infty)\), whose infimum is \(0\).

Problem (6C.17)

Suppose \(U\), \(V\), and \(W\) are normed vector spaces and \(T : U \to V\) and \(S : V \to W\) are linear. Prove that \(\|S \circ T\| \le \|S\| \, \|T\|\).

Solution

First note that \(S \circ T : U \to W\) is linear: if \(f, g \in U\) and \(\alpha \in \mathbf{F}\), then

\begin{equation*} (S \circ T)(f + g) = S\bigl(T(f+g)\bigr) = S(Tf + Tg) = S(Tf) + S(Tg) = (S \circ T)f + (S \circ T)g \end{equation*}

and

\begin{equation*} (S \circ T)(\alpha f) = S\bigl(T(\alpha f)\bigr) = S(\alpha \, Tf) = \alpha \, S(Tf) = \alpha \, (S \circ T)f . \end{equation*}

Hence \(\|S \circ T\|\) is defined by 6.43.

We use the inequality established in Exercise 6C.16(c): for any linear map \(R\) between normed vector spaces and any \(f\) in its domain,

\begin{equation*} \|Rf\| \le \|R\| \, \|f\| . \end{equation*}

(Proof, repeated for completeness: this is trivial for \(f = 0\) since \(R0 = 0\); and for \(f \ne 0\) the element \(u = f/\|f\|\) has \(\|u\| = 1 \le 1\), so \(\|Ru\| \le \|R\|\) by 6.43, and multiplying by \(\|f\|\) gives \(\|Rf\| = \|f\| \, \|Ru\| \le \|R\| \, \|f\|\).)

Case 1: \(\|S\| \ne 0\) and \(\|T\| \ne 0\). If either of \(\|S\|, \|T\|\) equals \(\infty\), then \(\|S\| \, \|T\| = \infty\) (both factors being nonzero) and the asserted inequality is immediate. So assume \(\|S\|, \|T\| \in (0, \infty)\). Let \(f \in U\) with \(\|f\| \le 1\). Applying the displayed inequality first to \(S\) (at the point \(Tf \in V\)) and then to \(T\) (at the point \(f\)),

\begin{equation*} \|(S \circ T)f\| = \|S(Tf)\| \le \|S\| \, \|Tf\| \le \|S\| \, \|T\| \, \|f\| \le \|S\| \, \|T\| , \end{equation*}

where the second inequality multiplies \(\|Tf\| \le \|T\| \, \|f\|\) by the nonnegative real number \(\|S\|\), and the last step uses \(\|f\| \le 1\). Taking the supremum over all \(f \in U\) with \(\|f\| \le 1\) gives

\begin{equation*} \|S \circ T\| \le \|S\| \, \|T\| . \end{equation*}

Case 2: \(\|S\| = 0\) or \(\|T\| = 0\). This case matters only because of the convention \(0 \cdot \infty = 0\) used throughout the book; we check that the inequality still holds. If \(\|T\| = 0\), then the displayed inequality gives \(\|Tf\| \le 0\), hence \(Tf = 0\), for every \(f \in U\); so \((S \circ T)f = S0 = 0\) for every \(f \in U\), and therefore \(\|S \circ T\| = 0 \le \|S\| \, \|T\|\). Similarly, if \(\|S\| = 0\), then \(Sg = 0\) for every \(g \in V\), so \((S \circ T)f = S(Tf) = 0\) for every \(f \in U\) and again \(\|S \circ T\| = 0 \le \|S\| \, \|T\|\).

In all cases \(\|S \circ T\| \le \|S\| \, \|T\|\), as desired. In particular, the composition of two bounded linear maps is bounded.

Problem (6C.18)

Suppose \(V\) and \(W\) are normed vector spaces and \(T : V \to W\) is a linear map. Prove that the following are equivalent.

(a) \(T\) is bounded.

(b) There exists \(f \in V\) such that \(T\) is continuous at \(f\).

(c) \(T\) is uniformly continuous (which means that for every \(\varepsilon > 0\), there exists \(\delta > 0\) such that \(\|Tf - Tg\| < \varepsilon\) for all \(f, g \in V\) with \(\|f - g\| < \delta\)).

(d) \(T^{-1}\bigl(B(0, r)\bigr)\) is an open subset of \(V\) for some \(r > 0\).

Solution

Here \(B(0,r) = \{h \in W : \|h\| < r\}\) is the open ball in \(W\) of radius \(r\) centered at \(0\), and \(T^{-1}\bigl(B(0,r)\bigr) = \{f \in V : \|Tf\| < r\}\).

We prove (a) \(\Rightarrow\) (c) \(\Rightarrow\) (b) \(\Rightarrow\) (a), and then (a) \(\Rightarrow\) (d) \(\Rightarrow\) (a); together these give all the equivalences.

Two facts will be used. First, \(T0 = 0\) and \(Tf - Tg = T(f - g)\) for all \(f, g \in V\), by linearity. Second, we record the following.

Lemma. Suppose \(\delta > 0\) and \(M \in [0, \infty)\) are such that \(\|Th\| \le M\) for all \(h \in V\) with \(\|h\| < \delta\). Then \(\|T\| \le 2M/\delta\); in particular \(T\) is bounded.

Proof of Lemma. Suppose \(f \in V\) with \(\|f\| \le 1\). Put \(h = (\delta/2) f\). Then \(\|h\| = (\delta/2)\|f\| \le \delta/2 < \delta\), so \(\|Th\| \le M\). But \(\|Th\| = (\delta/2)\|Tf\|\) by homogeneity of the norm on \(W\) and linearity of \(T\), so \(\|Tf\| \le 2M/\delta\). Taking the supremum over all \(f \in V\) with \(\|f\| \le 1\) gives \(\|T\| \le 2M/\delta < \infty\), as claimed.

(a) \(\Rightarrow\) (c). Suppose \(T\) is bounded, so \(\|T\| < \infty\). By Exercise 6C.16(c) we have \(\|Th\| \le \|T\| \, \|h\|\) for all \(h \in V\). Let \(\varepsilon > 0\) and set

\begin{equation*} \delta = \frac{\varepsilon}{1 + \|T\|} > 0 . \end{equation*}

If \(f, g \in V\) with \(\|f - g\| < \delta\), then, because \(\|T\| \ge 0\) and \(\|f - g\| < \delta\),

\begin{equation*} \|Tf - Tg\| = \|T(f - g)\| \le \|T\| \, \|f - g\| \le \|T\| \, \delta = \frac{\|T\|}{1 + \|T\|} \, \varepsilon < \varepsilon , \end{equation*}

the last inequality because \(\|T\| < 1 + \|T\|\) and \(\varepsilon > 0\). (Multiplying the strict inequality \(\|f - g\| < \delta\) by \(\|T\|\) yields only \(\le\), since \(\|T\|\) may be \(0\); this is why the chain above is written with \(\le\) at that step.) Thus \(T\) is uniformly continuous.

(c) \(\Rightarrow\) (b). Suppose \(T\) is uniformly continuous. Take \(f = 0 \in V\) (the set \(V\) is nonempty because it is a vector space). Given \(\varepsilon > 0\), uniform continuity supplies \(\delta > 0\) such that \(\|Tg - Th\| < \varepsilon\) whenever \(\|g - h\| < \delta\); in particular \(\|Tg - T0\| < \varepsilon\) for all \(g \in V\) with \(\|g - 0\| < \delta\). By the definition of continuity at a point (6.10), \(T\) is continuous at \(0\). So there exists an \(f \in V\) at which \(T\) is continuous.

(b) \(\Rightarrow\) (a). Suppose \(T\) is continuous at some \(f \in V\). Applying 6.10 with \(\varepsilon = 1\), there exists \(\delta > 0\) such that

\begin{equation*} \|Tg - Tf\| < 1 \quad \text{for all } g \in V \text{ with } \|g - f\| < \delta . \end{equation*}

Now suppose \(h \in V\) with \(\|h\| < \delta\). Then \(g = f + h\) satisfies \(\|g - f\| = \|h\| < \delta\), so

\begin{equation*} \|Th\| = \|T(g - f)\| = \|Tg - Tf\| < 1 . \end{equation*}

Thus \(\|Th\| \le 1\) whenever \(\|h\| < \delta\), and the Lemma (with \(M = 1\)) shows that \(\|T\| \le 2/\delta < \infty\). Hence \(T\) is bounded.

(a) \(\Rightarrow\) (d). Throughout this implication and the next, \(V\) and \(W\) are regarded as metric spaces with the metrics \(d(f,g) = \|f - g\|\) supplied by 6.36, so that the results of Section 6A apply. Suppose \(T\) is bounded. Then \(T\) is continuous, by 6.48 (or by the implication (a) \(\Rightarrow\) (c) already proved, since uniform continuity implies continuity). The ball \(B(0,1)\) is an open subset of \(W\) by 6.5, so \(T^{-1}\bigl(B(0,1)\bigr)\) is an open subset of \(V\) by the equivalence of 6.11(a) and 6.11(c). Thus (d) holds with \(r = 1\).

(d) \(\Rightarrow\) (a). Suppose \(r > 0\) and \(G = T^{-1}\bigl(B(0,r)\bigr)\) is an open subset of \(V\). Because \(T0 = 0\) and \(\|0\| = 0 < r\), we have \(0 \in G\). Since \(G\) is open, the definition of an open set (6.4) provides \(\delta > 0\) with \(B(0, \delta) \subseteq G\); that is,

\begin{equation*} h \in V \text{ and } \|h\| < \delta \implies \|Th\| < r . \end{equation*}

By the Lemma (with \(M = r\)), \(\|T\| \le 2r/\delta < \infty\), so \(T\) is bounded.

This completes the cycle (a) \(\Rightarrow\) (c) \(\Rightarrow\) (b) \(\Rightarrow\) (a), which shows (a), (b), (c) are equivalent, and the pair of implications (a) \(\Rightarrow\) (d) \(\Rightarrow\) (a), which adjoins (d). Hence all four statements are equivalent.

Note that (b) is a strikingly weak-looking hypothesis: continuity at a single point of \(V\) already forces uniform continuity on all of \(V\). This is a manifestation of the translation invariance built into linearity, which is also what makes 6.48 true.

Exercises 6D

Problem (6D.1)

Suppose \(V\) is a normed vector space and \(\varphi\) is a linear functional on \(V\). Suppose \(\alpha \in \mathbb{F} \setminus \{0\}\). Prove that the following are equivalent.

(a) \(\varphi\) is a bounded linear functional.

(b) \(\varphi^{-1}(\alpha)\) is a closed subset of \(V\).

(c) \(\overline{\varphi^{-1}(\alpha)} \neq V\).

Solution

Write \(S = \varphi^{-1}(\alpha) = \{f \in V : \varphi(f) = \alpha\}\).

We first record a fact used twice below. For \(h \in V\) the translation map \(\tau_h : V \to V\) defined by \(\tau_h(f) = h + f\) is a bijection of \(V\) satisfying

\begin{equation*} \|\tau_h(f) - \tau_h(g)\| = \|(h+f) - (h+g)\| = \|f - g\| \end{equation*}

for all \(f, g \in V\). Thus \(\tau_h\) and \(\tau_h^{-1} = \tau_{-h}\) are continuous, so \(\tau_h\) is a homeomorphism of \(V\) onto \(V\). In particular, for every \(A \subseteq V\) we have

\begin{equation*} \overline{h + A} = h + \overline{A}, \end{equation*}

and \(h + A\) is closed if and only if \(A\) is closed.

Case 1: \(\varphi\) is identically \(0\). Because \(\alpha \neq 0\), no \(f \in V\) satisfies \(\varphi(f) = \alpha\), so \(S = \varnothing\). Now (a) holds (\(\|\varphi\| = 0\)), (b) holds (\(\varnothing\) is closed), and (c) holds because \(\overline{\varnothing} = \varnothing \neq V\) (indeed \(0 \in V\), so \(V \neq \varnothing\)). All three statements are true, hence equivalent.

Case 2: \(\varphi\) is not identically \(0\). Then 6.52 applies to \(\varphi\). Choose \(g \in V\) with \(\varphi(g) = \beta \neq 0\) and set \(h = (\alpha/\beta) g\), so that \(\varphi(h) = \alpha\). For \(f \in V\),

\begin{equation*} \varphi(f) = \alpha \iff \varphi(f - h) = 0 \iff f - h \in \operatorname{null} \varphi , \end{equation*}

so

\begin{equation*} S = h + \operatorname{null} \varphi . \end{equation*}

(a) \(\Rightarrow\) (b). Suppose \(\varphi\) is bounded. By 6.48 the linear functional \(\varphi\) is continuous. The singleton \(\{\alpha\}\) is a closed subset of \(\mathbb{F}\), so \(S = \varphi^{-1}(\{\alpha\})\) is a closed subset of \(V\) by 6.11(d).

(b) \(\Rightarrow\) (c). Suppose \(S\) is closed, so \(\overline{S} = S\). Because \(\varphi(0) = 0 \neq \alpha\), we have \(0 \notin S\), and hence \(\overline{S} = S \neq V\).

(c) \(\Rightarrow\) (a). We prove the contrapositive. Suppose \(\varphi\) is not bounded. Then statement (a) of 6.52 fails for \(\varphi\), so statement (d) of 6.52 also fails; that is,

\begin{equation*} \overline{\operatorname{null} \varphi} = V . \end{equation*}

Using the translation fact with \(h\) as above,

\begin{equation*} \overline{S} = \overline{h + \operatorname{null} \varphi} = h + \overline{\operatorname{null} \varphi} = h + V = V . \end{equation*}

Thus (c) fails, completing the proof of the contrapositive.

Hence (a), (b), (c) are equivalent in both cases.

Problem (6D.2)

Suppose \(\varphi\) is a linear functional on a vector space \(V\). Prove that if \(U\) is a subspace of \(V\) such that \(\operatorname{null} \varphi \subseteq U\), then \(U = \operatorname{null} \varphi\) or \(U = V\).

Solution

Suppose \(U\) is a subspace of \(V\) with \(\operatorname{null} \varphi \subseteq U\). If \(U = \operatorname{null} \varphi\) we are done, so assume \(U \neq \operatorname{null} \varphi\). We must show \(U = V\).

Because \(\operatorname{null} \varphi \subseteq U\) and \(U \neq \operatorname{null} \varphi\), there exists \(u \in U\) with \(u \notin \operatorname{null} \varphi\). Set

\begin{equation*} \beta = \varphi(u) \neq 0 . \end{equation*}

Now let \(g \in V\) be arbitrary. Then

\begin{equation*} \varphi\Bigl(g - \frac{\varphi(g)}{\beta} u\Bigr) = \varphi(g) - \frac{\varphi(g)}{\beta}\,\varphi(u) = \varphi(g) - \varphi(g) = 0 , \end{equation*}

so \(g - \frac{\varphi(g)}{\beta} u \in \operatorname{null} \varphi \subseteq U\). Also \(\frac{\varphi(g)}{\beta} u \in U\) because \(U\) is a subspace containing \(u\). Because \(U\) is closed under addition,

\begin{equation*} g = \Bigl(g - \frac{\varphi(g)}{\beta} u\Bigr) + \frac{\varphi(g)}{\beta} u \in U . \end{equation*}

Because \(g \in V\) was arbitrary, \(V \subseteq U\). As \(U \subseteq V\), we conclude \(U = V\).

Thus \(U = \operatorname{null} \varphi\) or \(U = V\), as desired. (Note that the two alternatives are not exclusive: if \(\varphi\) is identically \(0\), then \(\operatorname{null} \varphi = V\) and the only subspace \(U\) containing \(\operatorname{null} \varphi\) is \(V\) itself, which equals both.)

Problem (6D.3)

Suppose \(\varphi\) and \(\psi\) are linear functionals on the same vector space. Prove that

\begin{equation*} \operatorname{null} \varphi \subseteq \operatorname{null} \psi \end{equation*}

if and only if there exists \(\alpha \in \mathbb{F}\) such that \(\psi = \alpha \varphi\).

Solution

Let \(V\) denote the vector space on which \(\varphi\) and \(\psi\) are defined.

First suppose there exists \(\alpha \in \mathbb{F}\) with \(\psi = \alpha \varphi\). If \(f \in \operatorname{null} \varphi\), then

\begin{equation*} \psi(f) = \alpha \varphi(f) = \alpha \cdot 0 = 0 , \end{equation*}

so \(f \in \operatorname{null} \psi\). Thus \(\operatorname{null} \varphi \subseteq \operatorname{null} \psi\), completing the proof in one direction.

To prove the other direction, suppose now that \(\operatorname{null} \varphi \subseteq \operatorname{null} \psi\).

If \(\varphi\) is identically \(0\), then \(\operatorname{null} \varphi = V\), so \(V = \operatorname{null} \varphi \subseteq \operatorname{null} \psi \subseteq V\), which forces \(\operatorname{null} \psi = V\); in other words \(\psi\) is identically \(0\). Hence \(\psi = 0 \cdot \varphi\), and we may take \(\alpha = 0\).

Now suppose \(\varphi\) is not identically \(0\). Choose \(g \in V\) with \(\varphi(g) \neq 0\) and set

\begin{equation*} h = \frac{g}{\varphi(g)} , \qquad \text{so that} \qquad \varphi(h) = 1 . \end{equation*}

Define

\begin{equation*} \alpha = \psi(h) . \end{equation*}

Let \(f \in V\). Then

\begin{equation*} \varphi\bigl(f - \varphi(f) h\bigr) = \varphi(f) - \varphi(f)\varphi(h) = \varphi(f) - \varphi(f) = 0 , \end{equation*}

so \(f - \varphi(f) h \in \operatorname{null} \varphi \subseteq \operatorname{null} \psi\). Therefore

\begin{equation*} 0 = \psi\bigl(f - \varphi(f) h\bigr) = \psi(f) - \varphi(f) \psi(h) = \psi(f) - \alpha \varphi(f) , \end{equation*}

where we used the linearity of \(\psi\) and the fact that \(\varphi(f)\) is a scalar. Hence \(\psi(f) = \alpha \varphi(f)\).

Because \(f \in V\) was arbitrary, \(\psi = \alpha \varphi\), completing the proof in the other direction.

Problem (6D.4)

[For this exercise and the next, \(\mathbb{F}^n\) should be endowed with the norm \(\|\cdot\|_\infty\) as defined in Example 6.34.]

Suppose \(n \in \mathbb{Z}^+\) and \(V\) is a normed vector space. Prove that every linear map from \(\mathbb{F}^n\) to \(V\) is continuous.

Solution

Recall from Example 6.34 that

\begin{equation*} \|(a_1, \ldots, a_n)\|_\infty = \max\{|a_1|, \ldots, |a_n|\} . \end{equation*}

Let \(T : \mathbb{F}^n \to V\) be a linear map. Let \(e_1, \ldots, e_n\) denote the standard basis of \(\mathbb{F}^n\), where \(e_k\) has \(1\) in the \(k^{\text{th}}\) slot and \(0\) elsewhere. Set

\begin{equation*} c = \sum_{k=1}^{n} \|T e_k\| \in [0, \infty) . \end{equation*}

Suppose \(x = (a_1, \ldots, a_n) \in \mathbb{F}^n\). Then \(x = \sum_{k=1}^{n} a_k e_k\), so the linearity of \(T\) and the triangle inequality and homogeneity of the norm on \(V\) give

\begin{equation*} \|T x\| = \Bigl\| \sum_{k=1}^{n} a_k \, T e_k \Bigr\| \leq \sum_{k=1}^{n} |a_k| \, \|T e_k\| \leq \Bigl( \max_{1 \leq k \leq n} |a_k| \Bigr) \sum_{k=1}^{n} \|T e_k\| = c \, \|x\|_\infty , \end{equation*}

where the second inequality holds because \(|a_k| \leq \|x\|_\infty\) for each \(k\).

Thus

\begin{equation*} \sup \{ \|Tx\| : x \in \mathbb{F}^n \text{ and } \|x\|_\infty \leq 1 \} \leq c < \infty , \end{equation*}

which means \(T\) is a bounded linear map, with \(\|T\| \leq c\). By 6.48, a linear map between normed vector spaces is continuous if and only if it is bounded. Hence \(T\) is continuous.

Problem (6D.5)

[For this exercise and the preceding one, \(\mathbb{F}^n\) should be endowed with the norm \(\|\cdot\|_\infty\) as defined in Example 6.34.]

Suppose \(n \in \mathbb{Z}^+\), \(V\) is a normed vector space, and \(T : \mathbb{F}^n \to V\) is a linear map that is one-to-one and onto \(V\).

(a) Show that

\begin{equation*} \inf \{ \|Tx\| : x \in \mathbb{F}^n \text{ and } \|x\|_\infty = 1 \} > 0 . \end{equation*}

(b) Prove that \(T^{-1} : V \to \mathbb{F}^n\) is a bounded linear map.

Solution

(a) Let

\begin{equation*} S = \{ x \in \mathbb{F}^n : \|x\|_\infty = 1 \} \quad\text{and}\quad c = \inf \{ \|Tx\| : x \in S \} . \end{equation*}

Note \(S \neq \varnothing\) (for example \(e_1 \in S\)) and \(c \geq 0\). Suppose, for contradiction, that \(c = 0\).

Then for each \(k \in \mathbb{Z}^+\) there exists \(x_k \in S\) with

\begin{equation*} \|T x_k\| < \tfrac{1}{k} . \end{equation*}

Write \(x_k = (a_{k,1}, \ldots, a_{k,n})\). For each \(j \in \{1, \ldots, n\}\) and each \(k\) we have \(|a_{k,j}| \leq \|x_k\|_\infty = 1\), so each of the \(n\) scalar sequences \((a_{k,j})_{k \in \mathbb{Z}^+}\) is a bounded sequence in \(\mathbb{F}\).

Apply the Bolzano–Weierstrass Theorem \(n\) times, passing successively to subsequences: first choose a subsequence along which the first coordinates converge, then a further subsequence along which the second coordinates converge, and so on. (When \(\mathbb{F} = \mathbb{C}\), a bounded sequence in \(\mathbb{C}\) has a convergent subsequence because its real and imaginary parts are bounded real sequences.) After \(n\) steps we obtain a subsequence \(x_{k_1}, x_{k_2}, \ldots\) and scalars \(a_1, \ldots, a_n\) such that

\begin{equation*} \lim_{m \to \infty} a_{k_m, j} = a_j \quad \text{for each } j \in \{1, \ldots, n\} . \end{equation*}

Put \(x = (a_1, \ldots, a_n)\). Because there are only finitely many coordinates,

\begin{equation*} \lim_{m \to \infty} \|x_{k_m} - x\|_\infty = \lim_{m \to \infty} \max_{1 \leq j \leq n} |a_{k_m, j} - a_j| = 0 . \end{equation*}

The reverse triangle inequality gives \(\bigl| \|x\|_\infty - \|x_{k_m}\|_\infty \bigr| \leq \|x - x_{k_m}\|_\infty \to 0\), and \(\|x_{k_m}\|_\infty = 1\) for every \(m\); hence \(\|x\|_\infty = 1\). In particular \(x \neq 0\).

By Exercise 6D.4, \(T\) is continuous, so \(T x_{k_m} \to Tx\) in \(V\), and therefore \(\|T x_{k_m}\| \to \|Tx\|\) (again by the reverse triangle inequality, this time in \(V\)). But \(\|T x_{k_m}\| < 1/k_m \to 0\), so \(\|Tx\| = 0\), that is, \(Tx = 0\).

Thus \(x\) is a nonzero element of \(\operatorname{null} T\), contradicting the hypothesis that \(T\) is one-to-one. Hence \(c > 0\), as desired.

(b) First we check that \(T^{-1}\) is linear. Because \(T\) is one-to-one and onto \(V\), the inverse function \(T^{-1} : V \to \mathbb{F}^n\) is well defined. Suppose \(g, h \in V\) and \(\alpha \in \mathbb{F}\). Then

\begin{equation*} T\bigl( T^{-1} g + \alpha \, T^{-1} h \bigr) = T(T^{-1} g) + \alpha \, T(T^{-1} h) = g + \alpha h , \end{equation*}

using the linearity of \(T\). Applying \(T^{-1}\) to both sides gives

\begin{equation*} T^{-1}(g + \alpha h) = T^{-1} g + \alpha \, T^{-1} h , \end{equation*}

so \(T^{-1}\) is linear.

Now let \(c > 0\) be the infimum from part (a). We claim

\begin{equation*} \|Tx\| \geq c \, \|x\|_\infty \quad \text{for all } x \in \mathbb{F}^n . \end{equation*}

Indeed, this is trivial if \(x = 0\) (both sides are \(0\)). If \(x \neq 0\), then \(\|x\|_\infty > 0\) and \(x / \|x\|_\infty \in S\), so by the definition of \(c\) and the homogeneity of the norm,

\begin{equation*} c \leq \Bigl\| T \Bigl( \frac{x}{\|x\|_\infty} \Bigr) \Bigr\| = \frac{\|Tx\|}{\|x\|_\infty} , \end{equation*}

which gives the claim.

Finally, suppose \(g \in V\) and set \(x = T^{-1} g\), so that \(Tx = g\). The claim gives

\begin{equation*} \|T^{-1} g\|_\infty = \|x\|_\infty \leq \frac{1}{c} \|Tx\| = \frac{1}{c} \|g\| . \end{equation*}

Hence

\begin{equation*} \sup \{ \|T^{-1} g\|_\infty : g \in V \text{ and } \|g\| \leq 1 \} \leq \frac{1}{c} < \infty , \end{equation*}

so \(T^{-1}\) is a bounded linear map, with \(\|T^{-1}\| \leq 1/c\).

Problem (6D.6)

Suppose \(n \in \mathbb{Z}^+\).

(a) Prove that all norms on \(\mathbb{F}^n\) have the same convergent sequences, the same open sets, and the same closed sets.

(b) Prove that all norms on \(\mathbb{F}^n\) make \(\mathbb{F}^n\) into a Banach space.

Solution

The key step is the following comparison with the norm \(\|\cdot\|_\infty\) of Example 6.34.

Claim. If \(\|\cdot\|\) is any norm on \(\mathbb{F}^n\), then there exist \(c, C \in (0, \infty)\) such that

\begin{equation*} c \, \|x\|_\infty \leq \|x\| \leq C \, \|x\|_\infty \quad \text{for all } x \in \mathbb{F}^n . \end{equation*}

Proof of claim. Let \(W\) denote the normed vector space \((\mathbb{F}^n, \|\cdot\|)\), and let \(I : \mathbb{F}^n \to W\) be the identity map \(I x = x\), where the domain carries \(\|\cdot\|_\infty\). Then \(I\) is linear, one-to-one, and onto \(W\).

Upper bound. Let \(e_1, \ldots, e_n\) be the standard basis of \(\mathbb{F}^n\) and set \(C = 1 + \sum_{k=1}^{n} \|e_k\|\), where the norm is \(\|\cdot\|\). Then \(C \in (0,\infty)\), and for \(x = (a_1, \ldots, a_n) \in \mathbb{F}^n\) the triangle inequality and homogeneity of \(\|\cdot\|\) give

\begin{equation*} \|x\| = \Bigl\| \sum_{k=1}^{n} a_k e_k \Bigr\| \leq \sum_{k=1}^{n} |a_k| \, \|e_k\| \leq \Bigl( \max_{1 \leq k \leq n} |a_k| \Bigr) \sum_{k=1}^{n} \|e_k\| \leq C \, \|x\|_\infty . \end{equation*}

(This is exactly the computation in Exercise 6D.4 applied to \(I\); in particular \(I\) is a bounded, hence continuous, linear map from \((\mathbb{F}^n, \|\cdot\|_\infty)\) to \(W\).)

Lower bound. Because \(I\) is linear, one-to-one, and onto the normed vector space \(W\), Exercise 6D.5(a) applies to \(I\) and yields

\begin{equation*} c := \inf \{ \|Ix\| : x \in \mathbb{F}^n \text{ and } \|x\|_\infty = 1 \} = \inf \{ \|x\| : \|x\|_\infty = 1 \} > 0 . \end{equation*}

Now let \(x \in \mathbb{F}^n\). If \(x = 0\) then \(c \|x\|_\infty = 0 = \|x\|\). If \(x \neq 0\) then \(\|x\|_\infty > 0\) and \(x / \|x\|_\infty\) has \(\|\cdot\|_\infty\)-norm \(1\), so by the definition of \(c\) and homogeneity of \(\|\cdot\|\),

\begin{equation*} c \leq \Bigl\| \frac{x}{\|x\|_\infty} \Bigr\| = \frac{\|x\|}{\|x\|_\infty} , \end{equation*}

that is, \(c \|x\|_\infty \leq \|x\|\). This proves the claim.

Consequently, if \(\|\cdot\|_a\) and \(\|\cdot\|_b\) are any two norms on \(\mathbb{F}^n\), then applying the claim to each and combining, there exist \(\lambda, \Lambda \in (0,\infty)\) with

\begin{equation*} \lambda \, \|x\|_b \leq \|x\|_a \leq \Lambda \, \|x\|_b \quad \text{for all } x \in \mathbb{F}^n . \end{equation*}

Explicitly, if \(c_a \|x\|_\infty \leq \|x\|_a \leq C_a \|x\|_\infty\) and \(c_b \|x\|_\infty \leq \|x\|_b \leq C_b \|x\|_\infty\), then \(\|x\|_a \leq C_a \|x\|_\infty \leq (C_a / c_b) \|x\|_b\) and \(\|x\|_a \geq c_a \|x\|_\infty \geq (c_a / C_b) \|x\|_b\); take \(\Lambda = C_a/c_b\) and \(\lambda = c_a/C_b\).

(a) Fix two norms \(\|\cdot\|_a\), \(\|\cdot\|_b\) on \(\mathbb{F}^n\) and let \(\lambda, \Lambda \in (0,\infty)\) be as above.

Convergent sequences. Suppose \(x_1, x_2, \ldots\) is a sequence in \(\mathbb{F}^n\) and \(x \in \mathbb{F}^n\). Because

\begin{equation*} \lambda \, \|x_k - x\|_b \leq \|x_k - x\|_a \leq \Lambda \, \|x_k - x\|_b , \end{equation*}

we have \(\|x_k - x\|_a \to 0\) if and only if \(\|x_k - x\|_b \to 0\). Thus the two norms have exactly the same convergent sequences, with the same limits.

Open sets. For \(y \in \mathbb{F}^n\) and \(r > 0\) let \(B_a(y, r)\) and \(B_b(y, r)\) denote the corresponding open balls. If \(\|z - y\|_b < r/\Lambda\), then \(\|z - y\|_a \leq \Lambda \|z - y\|_b < r\); hence

\begin{equation*} B_b(y, r/\Lambda) \subseteq B_a(y, r) . \end{equation*}

Now suppose \(G \subseteq \mathbb{F}^n\) is open with respect to \(\|\cdot\|_a\) and \(y \in G\). Choose \(r > 0\) with \(B_a(y,r) \subseteq G\). Then \(B_b(y, r/\Lambda) \subseteq G\). Thus \(G\) is open with respect to \(\|\cdot\|_b\). Interchanging the roles of the two norms (using \(\|z-y\|_b \leq \lambda^{-1}\|z-y\|_a\)) gives the reverse implication. Hence the two norms have the same open sets.

Closed sets. A subset of \(\mathbb{F}^n\) is closed with respect to a norm if and only if its complement is open with respect to that norm. Because the two norms have the same open sets, they have the same closed sets.

Because \(\|\cdot\|_a\) and \(\|\cdot\|_b\) were arbitrary norms on \(\mathbb{F}^n\), part (a) is proved.

(b) Let \(\|\cdot\|\) be a norm on \(\mathbb{F}^n\) and take \(c, C \in (0,\infty)\) as in the claim. We must show that \((\mathbb{F}^n, \|\cdot\|)\) is complete.

Suppose \(x_1, x_2, \ldots\) is a Cauchy sequence with respect to \(\|\cdot\|\), and write \(x_k = (a_{k,1}, \ldots, a_{k,n})\). For all \(j, k\),

\begin{equation*} \|x_j - x_k\|_\infty \leq \tfrac{1}{c} \|x_j - x_k\| , \end{equation*}

so the sequence is also Cauchy with respect to \(\|\cdot\|_\infty\). For each fixed \(m \in \{1, \ldots, n\}\) we have

\begin{equation*} |a_{j,m} - a_{k,m}| \leq \|x_j - x_k\|_\infty , \end{equation*}

so \((a_{k,m})_{k \in \mathbb{Z}^+}\) is a Cauchy sequence in \(\mathbb{F}\). Because \(\mathbb{F}\) (which is \(\mathbb{R}\) or \(\mathbb{C}\)) is complete, this sequence converges; call its limit \(a_m\). Put \(x = (a_1, \ldots, a_n) \in \mathbb{F}^n\).

Because the maximum is over finitely many indices,

\begin{equation*} \|x_k - x\|_\infty = \max_{1 \leq m \leq n} |a_{k,m} - a_m| \longrightarrow 0 \quad \text{as } k \to \infty . \end{equation*}

Hence

\begin{equation*} \|x_k - x\| \leq C \, \|x_k - x\|_\infty \longrightarrow 0 , \end{equation*}

so \(x_k \to x\) with respect to \(\|\cdot\|\).

Thus every Cauchy sequence in \((\mathbb{F}^n, \|\cdot\|)\) converges in \((\mathbb{F}^n, \|\cdot\|)\), which means \((\mathbb{F}^n, \|\cdot\|)\) is a Banach space.

Problem (6D.7)

Suppose \(V\) and \(W\) are normed vector spaces and \(V\) is finite-dimensional. Prove that every linear map from \(V\) to \(W\) is continuous.

Solution

Let \(T : V \to W\) be a linear map. By 6.48 it suffices to show that \(T\) is bounded.

If \(V = \{0\}\), then \(T\) is the zero map and \(\|T\| = 0\), so \(T\) is bounded and hence continuous. So assume \(V \neq \{0\}\) and let \(n = \dim V \in \mathbb{Z}^+\). Choose a basis \(e_1, \ldots, e_n\) of \(V\) and define \(S : \mathbb{F}^n \to V\) by

\begin{equation*} S(a_1, \ldots, a_n) = a_1 e_1 + \cdots + a_n e_n , \end{equation*}

where \(\mathbb{F}^n\) carries the norm \(\|\cdot\|_\infty\) of Example 6.34. Then \(S\) is linear. Because \(e_1, \ldots, e_n\) spans \(V\), the map \(S\) is onto \(V\); because \(e_1, \ldots, e_n\) is linearly independent, \(S(a_1,\ldots,a_n) = 0\) forces \(a_1 = \cdots = a_n = 0\), so \(\operatorname{null} S = \{0\}\) and \(S\) is one-to-one.

We record the standard inequality that will be used twice: if \(R : X \to Y\) is a bounded linear map between normed vector spaces, then

\begin{equation*} \|R f\| \leq \|R\| \, \|f\| \quad \text{for all } f \in X . \end{equation*}

Indeed this is trivial when \(f = 0\), and when \(f \neq 0\) the vector \(f / \|f\|\) has norm \(1\), so \(\|R(f/\|f\|)\| \leq \|R\|\) by the definition of \(\|R\|\) in 6.43; multiplying by \(\|f\|\) and using homogeneity gives the inequality.

By Exercise 6D.5(b) (applicable because \(V\) is a normed vector space and \(S : \mathbb{F}^n \to V\) is linear, one-to-one, and onto \(V\)), the linear map \(S^{-1} : V \to \mathbb{F}^n\) is bounded. Set \(M = \|S^{-1}\|\), so that

\begin{equation*} \|S^{-1} f\|_\infty \leq M \|f\| \quad \text{for all } f \in V . \end{equation*}

Next consider \(T \circ S : \mathbb{F}^n \to W\). It is a composition of linear maps, hence linear, so by Exercise 6D.4 it is continuous and therefore bounded (6.48). Set \(K = \|T \circ S\|\), so that

\begin{equation*} \|(T \circ S)(x)\| \leq K \|x\|_\infty \quad \text{for all } x \in \mathbb{F}^n . \end{equation*}

Now let \(f \in V\). Taking \(x = S^{-1} f\) and using \(S(S^{-1} f) = f\), we get

\begin{equation*} \|T f\| = \|(T \circ S)(S^{-1} f)\| \leq K \, \|S^{-1} f\|_\infty \leq K M \, \|f\| . \end{equation*}

Hence

\begin{equation*} \sup \{ \|Tf\| : f \in V \text{ and } \|f\| \leq 1 \} \leq K M < \infty , \end{equation*}

so \(T\) is a bounded linear map, with \(\|T\| \leq KM\). By 6.48, \(T\) is continuous.

Because \(T\) was an arbitrary linear map from \(V\) to \(W\), every linear map from \(V\) to \(W\) is continuous.

Problem (6D.8)

Prove that every finite-dimensional normed vector space is a Banach space.

Solution

Suppose \(V\) is a finite-dimensional normed vector space. If \(V = \{0\}\) then \(V\) is trivially complete, so assume \(V \ne \{0\}\). By 6.54 there is a finite family spanning \(V\); discarding in turn each element that lies in the span of the others leaves a linearly independent finite spanning family, so \(V\) has a basis \(e_1, \ldots, e_n\) for some \(n \in \mathbf{Z}^+\). Throughout, \(\mathbf{F}^n\) is endowed with the norm \(\|\cdot\|_\infty\) of Example 6.34, so that

\begin{equation*} \|(a_1, \ldots, a_n)\|_\infty = \max\{|a_1|, \ldots, |a_n|\}. \end{equation*}

Define \(T \colon \mathbf{F}^n \to V\) by

\begin{equation*} T(a_1, \ldots, a_n) = a_1 e_1 + \cdots + a_n e_n . \end{equation*}

Then \(T\) is linear, and \(T\) is one-to-one and onto \(V\) because \(e_1, \ldots, e_n\) is a basis of \(V\).

Step 1: \(T\) is bounded. Let \(M = \|e_1\| + \cdots + \|e_n\|\). For \(x = (a_1, \ldots, a_n) \in \mathbf{F}^n\) the triangle inequality and homogeneity of the norm give

\begin{equation*} \|Tx\| \le \sum_{k=1}^{n} |a_k| \, \|e_k\| \le M \|x\|_\infty . \end{equation*}

Thus \(\|T\| \le M < \infty\).

Step 2: there exists \(c > 0\) such that \(\|Tx\| \ge c \|x\|_\infty\) for all \(x \in \mathbf{F}^n\). Suppose no such \(c\) exists. Then for each \(j \in \mathbf{Z}^+\) there is \(y_j \in \mathbf{F}^n\) with \(\|T y_j\| < \frac{1}{j} \|y_j\|_\infty\); in particular \(y_j \ne 0\), so setting \(x_j = y_j / \|y_j\|_\infty\) we obtain a sequence \(x_1, x_2, \ldots\) in \(\mathbf{F}^n\) with

\begin{equation*} \|x_j\|_\infty = 1 \quad \text{and} \quad \lim_{j \to \infty} \|T x_j\| = 0 . \end{equation*}

Each coordinate sequence of \((x_j)\) is a bounded sequence in \(\mathbf{F}\) (bounded by \(1\) in absolute value). Applying the Bolzano-Weierstrass Theorem successively in each of the \(n\) coordinates (in the complex case, apply it to real and imaginary parts) and passing to a subsequence \(n\) times, we obtain a subsequence \(x_{j_1}, x_{j_2}, \ldots\) that converges coordinatewise to some \(x \in \mathbf{F}^n\). Coordinatewise convergence of finitely many coordinates is exactly convergence in \(\|\cdot\|_\infty\), so \(\|x_{j_m} - x\|_\infty \to 0\) as \(m \to \infty\). Because \(\big| \|x_{j_m}\|_\infty - \|x\|_\infty \big| \le \|x_{j_m} - x\|_\infty\), we get \(\|x\|_\infty = 1\); in particular \(x \ne 0\).

Now Step 1 gives

\begin{equation*} \|Tx\| \le \|T(x - x_{j_m})\| + \|T x_{j_m}\| \le M \|x - x_{j_m}\|_\infty + \|T x_{j_m}\| , \end{equation*}

and both terms on the right tend to \(0\) as \(m \to \infty\). Hence \(Tx = 0\). Because \(T\) is one-to-one this forces \(x = 0\), contradicting \(\|x\|_\infty = 1\). This contradiction proves that the desired \(c > 0\) exists.

Step 3: \(V\) is complete. Suppose \(f_1, f_2, \ldots\) is a Cauchy sequence in \(V\). Because \(T\) is onto \(V\) and one-to-one, for each \(j\) there is a unique \(x_j \in \mathbf{F}^n\) with \(T x_j = f_j\). Step 2 applied to \(x_j - x_k\) gives

\begin{equation*} \|x_j - x_k\|_\infty \le \frac{1}{c} \|T x_j - T x_k\| = \frac{1}{c} \|f_j - f_k\| , \end{equation*}

so \(x_1, x_2, \ldots\) is a Cauchy sequence in \(\mathbf{F}^n\) with respect to \(\|\cdot\|_\infty\). Each coordinate sequence is then a Cauchy sequence in \(\mathbf{F}\), which is complete; hence the coordinate sequences converge, and therefore \(\|x_j - x\|_\infty \to 0\) for some \(x \in \mathbf{F}^n\). Setting \(f = Tx \in V\), Step 1 gives

\begin{equation*} \|f_j - f\| = \|T(x_j - x)\| \le M \|x_j - x\|_\infty \to 0 . \end{equation*}

Thus \(f_j \to f\) in \(V\).

Every Cauchy sequence in \(V\) converges in \(V\), so \(V\) is a Banach space.

Problem (6D.9)

Prove that every finite-dimensional subspace of each normed vector space is closed.

Solution

Suppose \(V\) is a normed vector space and \(U\) is a finite-dimensional subspace of \(V\). Give \(U\) the norm it inherits from \(V\); then \(U\) is a finite-dimensional normed vector space, so Exercise 8 in this section shows that \(U\) is a Banach space, meaning that every Cauchy sequence of elements of \(U\) converges to an element of \(U\).

Now suppose \(f \in \overline{U}\). Because the norm makes \(V\) into a metric space (with \(d(f, g) = \|f - g\|\)), there exists a sequence \(f_1, f_2, \ldots\) in \(U\) with \(\lim_{j \to \infty} \|f_j - f\| = 0\); indeed, for each \(j \in \mathbf{Z}^+\) the open ball of radius \(\frac{1}{j}\) centered at \(f\) intersects \(U\) (otherwise its complement would be a closed set containing \(U\) but not \(f\), contradicting \(f \in \overline{U}\)), and we may take \(f_j\) in that intersection.

A convergent sequence is a Cauchy sequence: given \(\varepsilon > 0\), choose \(J\) so that \(\|f_j - f\| < \frac{\varepsilon}{2}\) for all \(j \ge J\); then \(\|f_j - f_k\| \le \|f_j - f\| + \|f - f_k\| < \varepsilon\) for all \(j, k \ge J\). Thus \(f_1, f_2, \ldots\) is a Cauchy sequence of elements of \(U\), and hence, because \(U\) is a Banach space, there exists \(g \in U\) with \(\lim_{j \to \infty} \|f_j - g\| = 0\).

Limits in the metric space \(V\) are unique: \(\|f - g\| \le \|f - f_j\| + \|f_j - g\| \to 0\), so \(\|f - g\| = 0\) and therefore \(f = g\). Hence \(f \in U\).

We have shown \(\overline{U} \subseteq U\). Because \(U \subseteq \overline{U}\) always holds, \(\overline{U} = U\), which means that \(U\) is a closed subset of \(V\).

Problem (6D.10)

Give a concrete example of an infinite-dimensional normed vector space and a basis of that normed vector space.

Solution

Let

\begin{equation*} V = \{(a_1, a_2, \ldots) \in \ell^\infty : a_k = 0 \text{ for all but finitely many } k \in \mathbf{Z}^+\} , \end{equation*}

with the norm \(\|\cdot\|_\infty\) inherited from \(\ell^\infty\), so that \(\|(a_1, a_2, \ldots)\|_\infty = \sup_{k \in \mathbf{Z}^+} |a_k|\) (which is a maximum over a finite set here). For \(k \in \mathbf{Z}^+\) let \(e_k \in V\) be the sequence whose \(k^{\text{th}}\) coordinate is \(1\) and whose other coordinates are all \(0\). We claim that \(V\) is an infinite-dimensional normed vector space and that \(\{e_k\}_{k \in \mathbf{Z}^+}\) is a basis of \(V\), in the sense of 6.54 (only finite sums are involved).

\(V\) is a normed vector space: \(V\) is a subspace of \(\ell^\infty\), because \(0 \in V\) and because if \(a, b \in V\) and \(\alpha \in \mathbf{F}\), then the coordinates of \(a + b\) and of \(\alpha a\) vanish outside the union of the two finite sets on which the coordinates of \(a\) and of \(b\) are nonzero, which is again a finite set. The restriction of a norm to a subspace is a norm.

\(\{e_k\}_{k \in \mathbf{Z}^+}\) spans \(V\): if \(a = (a_1, a_2, \ldots) \in V\), choose \(n \in \mathbf{Z}^+\) such that \(a_k = 0\) for all \(k > n\); then \(a = \sum_{k=1}^{n} a_k e_k\), a finite linear combination of the \(e_k\).

\(\{e_k\}_{k \in \mathbf{Z}^+}\) is linearly independent: if \(\Omega \subseteq \mathbf{Z}^+\) is finite and \(\sum_{k \in \Omega} \alpha_k e_k = 0\), then for each \(j \in \Omega\) the \(j^{\text{th}}\) coordinate of the left side is \(\alpha_j\) and the \(j^{\text{th}}\) coordinate of the right side is \(0\); hence \(\alpha_j = 0\) for every \(j \in \Omega\).

Thus \(\{e_k\}_{k \in \mathbf{Z}^+}\) is a basis of \(V\).

\(V\) is infinite-dimensional: suppose \(V\) were spanned by a finite list \(f_1, \ldots, f_m\) of elements of \(V\). Each \(f_j\) has only finitely many nonzero coordinates, so there exists \(N \in \mathbf{Z}^+\) such that the \(k^{\text{th}}\) coordinate of every \(f_j\) equals \(0\) for all \(k > N\). Every element of \(\operatorname{span}(f_1, \ldots, f_m)\) then also has \(k^{\text{th}}\) coordinate \(0\) for all \(k > N\). But \(e_{N+1} \in V\) has \((N+1)^{\text{st}}\) coordinate equal to \(1\), so \(e_{N+1} \notin \operatorname{span}(f_1, \ldots, f_m)\), a contradiction. Hence \(V\) is not finite-dimensional.

Note that \(V\) is not a Banach space. For example, let \(b^{(n)} = \sum_{k=1}^{n} \frac{1}{k} e_k \in V\). If \(n > m\) then \(\|b^{(n)} - b^{(m)}\|_\infty = \max_{m < k \le n} \frac{1}{k} = \frac{1}{m + 1}\), so \(b^{(1)}, b^{(2)}, \ldots\) is a Cauchy sequence in \(V\). If \(b \in V\) satisfied \(\|b^{(n)} - b\|_\infty \to 0\), then, because convergence in \(\|\cdot\|_\infty\) implies convergence in each coordinate, the \(k^{\text{th}}\) coordinate of \(b\) would equal \(\frac{1}{k}\) for every \(k \in \mathbf{Z}^+\), contradicting that elements of \(V\) have only finitely many nonzero coordinates. The incompleteness is not an accident of this particular example: a Baire category argument shows that an infinite-dimensional Banach space never has a countable basis, so no example whose basis can be listed as \(e_1, e_2, \ldots\) could have been complete. As the text remarks, no one has ever produced a concrete example of a basis of an infinite-dimensional Banach space.

Another concrete example: let \(\mathcal{P}\) denote the vector space of polynomial functions from \([0, 1]\) to \(\mathbf{F}\), with the norm \(\|p\| = \sup_{[0,1]} |p|\). Then \(\{x^n\}_{n \in \{0, 1, 2, \ldots\}}\) is a basis of \(\mathcal{P}\): it spans \(\mathcal{P}\) by the definition of a polynomial, and it is linearly independent because a nonzero polynomial of degree \(n\) has at most \(n\) zeros and hence cannot be the zero function on \([0, 1]\). The argument above (comparing degrees instead of coordinates) shows that \(\mathcal{P}\) is infinite-dimensional.

Problem (6D.11)

Show that the collection \(\mathcal{A} = \{k\mathbf{Z} : k = 2, 3, 4, \ldots\}\) of subsets of \(\mathbf{Z}\) satisfies the hypothesis of Zorn’s Lemma (6.60).

Solution

Here \(k\mathbf{Z} = \{km : m \in \mathbf{Z}\}\), as in Example 6.56. The hypothesis of Zorn’s Lemma (6.60) with \(V = \mathbf{Z}\) is that the union of all the sets in \(\mathcal{C}\) belongs to \(\mathcal{A}\) for every chain \(\mathcal{C} \subseteq \mathcal{A}\).

First we record the divisibility criterion for inclusion. Suppose \(j, k \in \{2, 3, 4, \ldots\}\). Then

\begin{equation*} k\mathbf{Z} \subseteq j\mathbf{Z} \iff j \text{ divides } k . \end{equation*}

Indeed, if \(k\mathbf{Z} \subseteq j\mathbf{Z}\) then \(k = k \cdot 1 \in k\mathbf{Z} \subseteq j\mathbf{Z}\), which means \(k = jm\) for some \(m \in \mathbf{Z}\), so \(j\) divides \(k\). Conversely, if \(k = jm\) for some \(m \in \mathbf{Z}\), then every element \(kn\) of \(k\mathbf{Z}\) equals \(j(mn) \in j\mathbf{Z}\).

Now suppose \(\mathcal{C} \subseteq \mathcal{A}\) is a nonempty chain, and let

\begin{equation*} U = \bigcup_{E \in \mathcal{C}} E . \end{equation*}

Let \(S = \{k \in \{2, 3, 4, \ldots\} : k\mathbf{Z} \in \mathcal{C}\}\). Because \(\mathcal{C}\) is a nonempty subcollection of \(\mathcal{A}\), the set \(S\) is a nonempty set of integers each of which is at least \(2\); by the well-ordering property of the positive integers, \(S\) has a smallest element \(m\). Note \(U = \bigcup_{k \in S} k\mathbf{Z}\).

We claim that \(U = m\mathbf{Z}\).

Because \(m \in S\), we have \(m\mathbf{Z} \in \mathcal{C}\) and thus \(m\mathbf{Z} \subseteq U\).

For the other inclusion, suppose \(k \in S\). Then \(k\mathbf{Z}\) and \(m\mathbf{Z}\) both belong to the chain \(\mathcal{C}\), so \(k\mathbf{Z} \subseteq m\mathbf{Z}\) or \(m\mathbf{Z} \subseteq k\mathbf{Z}\). In the second case the divisibility criterion gives that \(k\) divides \(m\); because \(k \ge 2\) and \(m \ge 2\) are positive, this implies \(k \le m\), and the minimality of \(m\) in \(S\) then gives \(k = m\), so that \(k\mathbf{Z} \subseteq m\mathbf{Z}\) holds in this case too. Hence \(k\mathbf{Z} \subseteq m\mathbf{Z}\) for every \(k \in S\), and therefore \(U \subseteq m\mathbf{Z}\).

Thus \(U = m\mathbf{Z} \in \mathcal{A}\), as desired. (In fact the computation shows more: the union of a nonempty chain in \(\mathcal{A}\) is simply the largest set in that chain, which is \(m\mathbf{Z}\) for the smallest index \(m\) occurring in the chain.)

Only the empty chain remains. If \(\mathcal{C} = \emptyset\), then the union of all the sets in \(\mathcal{C}\) is \(\emptyset\), which is not an element of \(\mathcal{A}\). As is standard, chains here are understood to be nonempty (this is the only reading under which the hypothesis of 6.60 can be met by a collection that does not contain the empty set). Alternatively, one may replace \(\mathcal{A}\) by \(\mathcal{A} \cup \{\emptyset\}\): unions of nonempty chains are unchanged, unions of chains consisting only of \(\emptyset\) equal \(\emptyset\), and since \(\emptyset \subsetneq 2\mathbf{Z}\) the set \(\emptyset\) is not maximal, so the maximal elements of \(\mathcal{A} \cup \{\emptyset\}\) are exactly those of \(\mathcal{A}\).

Consequently Zorn’s Lemma (6.60) applies to \(\mathcal{A}\) and guarantees that \(\mathcal{A}\) has a maximal element, in agreement with Example 6.56, which identifies the maximal elements of \(\mathcal{A}\) as the sets \(p\mathbf{Z}\) with \(p\) prime.

Problem (6D.12)

Prove that every linearly independent family in a vector space can be extended to a basis of the vector space.

Solution

Suppose \(V\) is a vector space and \(\{f_j\}_{j \in \Lambda}\) is a linearly independent family in \(V\). First note that the map \(j \mapsto f_j\) is one-to-one: if \(i, j \in \Lambda\) with \(i \ne j\) and \(f_i = f_j\), then \(1 \cdot f_i + (-1) \cdot f_j = 0\) is a linear combination with coefficients not all \(0\) that equals \(0\), contradicting linear independence. Hence, as in the convention preceding 6.57, we may identify the family with the set

\begin{equation*} A = \{f_j : j \in \Lambda\} , \end{equation*}

which is a linearly independent subset of \(V\); extending the family to a basis means finding a basis of \(V\) that contains \(A\).

Let \(\mathcal{A}\) be the collection of all linearly independent subsets \(E\) of \(V\) such that \(A \subseteq E\). Note \(A \in \mathcal{A}\), so \(\mathcal{A} \ne \emptyset\).

We verify the hypothesis of Zorn’s Lemma (6.60). Suppose \(\mathcal{C} \subseteq \mathcal{A}\) is a nonempty chain, and let

\begin{equation*} U = \bigcup_{E \in \mathcal{C}} E . \end{equation*}

Because \(\mathcal{C}\) is nonempty and every element of \(\mathcal{C}\) contains \(A\), we have \(A \subseteq U\).

To see that \(U\) is linearly independent, first observe that every finite subset of \(U\) is contained in a single element of \(\mathcal{C}\). To prove this, suppose \(u_1, \ldots, u_n \in U\) and choose \(E_1, \ldots, E_n \in \mathcal{C}\) with \(u_k \in E_k\) for each \(k\). We show by induction on \(n\) that one of \(E_1, \ldots, E_n\) contains all the others. This is trivial for \(n = 1\). If it holds for \(n - 1\), say \(E_i \supseteq E_k\) for all \(k \le n - 1\), then, because \(\mathcal{C}\) is a chain, either \(E_n \subseteq E_i\) (and \(E_i\) contains all of \(E_1, \ldots, E_n\)) or \(E_i \subseteq E_n\) (and then \(E_n\) contains all of \(E_1, \ldots, E_n\)). Thus some \(E_i\) contains \(u_1, \ldots, u_n\).

Linear independence is a condition that involves only finitely many elements at a time: if \(\Omega\) is a finite subset of \(U\) and \(\{\alpha_u\}_{u \in \Omega}\) is a family in \(\mathbf{F}\) with \(\sum_{u \in \Omega} \alpha_u u = 0\), then by the previous paragraph \(\Omega \subseteq E\) for some \(E \in \mathcal{C}\); since \(E\) is linearly independent, \(\alpha_u = 0\) for every \(u \in \Omega\). Hence \(U\) is linearly independent, and therefore \(U \in \mathcal{A}\).

(As in Exercise 11, chains are understood to be nonempty. If \(A = \emptyset\), then even the empty chain causes no trouble, because its union \(\emptyset\) belongs to \(\mathcal{A}\) in that case. If \(A \ne \emptyset\), one may instead apply Zorn’s Lemma to \(\mathcal{A} \cup \{\emptyset\}\): the union of a chain in \(\mathcal{A} \cup \{\emptyset\}\) is either \(\emptyset\), which lies in \(\mathcal{A} \cup \{\emptyset\}\), or the union of the nonempty sets in that chain, which lies in \(\mathcal{A}\) by the paragraphs above; and \(\emptyset \subsetneq A\) shows \(\emptyset\) is not maximal, so the maximal elements of \(\mathcal{A} \cup \{\emptyset\}\) are exactly those of \(\mathcal{A}\).)

By Zorn’s Lemma (6.60), \(\mathcal{A}\) contains a maximal element \(\Gamma\). Thus \(\Gamma\) is a linearly independent subset of \(V\) with \(A \subseteq \Gamma\), and there is no linearly independent set \(\Gamma^{\prime}\) containing \(A\) with \(\Gamma \subsetneq \Gamma^{\prime}\).

We claim that \(\Gamma\) is in fact a maximal element of the collection of all linearly independent subsets of \(V\). Suppose \(\Gamma^{\prime}\) is a linearly independent subset of \(V\) with \(\Gamma \subsetneq \Gamma^{\prime}\). Then \(A \subseteq \Gamma \subseteq \Gamma^{\prime}\), so \(\Gamma^{\prime} \in \mathcal{A}\), contradicting the maximality of \(\Gamma\) in \(\mathcal{A}\). Hence no such \(\Gamma^{\prime}\) exists.

Now 6.57 implies that \(\Gamma\) is a basis of \(V\). Because \(A \subseteq \Gamma\), this basis extends the original linearly independent family \(\{f_j\}_{j \in \Lambda}\): writing \(\Gamma\) as the family \(\{e_f\}_{f \in \Gamma}\) with \(e_f = f\), the subfamily indexed by \(A\) is precisely the given family. Thus every linearly independent family in a vector space can be extended to a basis of that vector space.

Problem (6D.13)

Suppose \(V\) is a normed vector space, \(U\) is a subspace of \(V\), and \(\psi \colon U \to \mathbf{R}\) is a bounded linear functional. Prove that \(\psi\) has a unique extension to a bounded linear functional \(\varphi\) on \(V\) with \(\|\varphi\| = \|\psi\|\) if and only if

\begin{equation*} \sup_{f \in U} \big( -\|\psi\| \, \|f + h\| - \psi(f) \big) = \inf_{g \in U} \big( \|\psi\| \, \|g + h\| - \psi(g) \big) \end{equation*}

for every \(h \in V \setminus U\).

Solution

Because \(\psi\) maps into \(\mathbf{R}\), here \(\mathbf{F} = \mathbf{R}\) and \(V\) is a real normed vector space. (This is the only reading that is not degenerate: if \(\mathbf{F} = \mathbf{C}\), then \(U\) is closed under multiplication by \(i\), so a \(\mathbf{C}\)-linear \(\psi\) with real values would satisfy \(i \psi(f) = \psi(i f) \in \mathbf{R}\) for all \(f \in U\), forcing \(\psi = 0\).) For \(h \in V \setminus U\) write

\begin{equation*} s(h) = \sup_{f \in U} \big( -\|\psi\| \, \|f + h\| - \psi(f) \big), \qquad t(h) = \inf_{g \in U} \big( \|\psi\| \, \|g + h\| - \psi(g) \big) . \end{equation*}

Taking \(f = 0\) and \(g = 0\) shows that \(s(h) \ge -\|\psi\| \, \|h\|\) and \(t(h) \le \|\psi\| \, \|h\|\), and 6.66 (proved in the course of proving the Extension Lemma 6.63) states that

\begin{equation*} s(h) \le t(h) . \end{equation*}

Hence \(s(h)\) and \(t(h)\) are real numbers, and the displayed condition of the exercise says precisely that \(s(h) = t(h)\).

The following observation is the bridge between the two sides of the equivalence.

Claim. Suppose \(h \in V \setminus U\) and \(c \in \mathbf{R}\), and define \(\varphi_c \colon U + \mathbf{R}h \to \mathbf{R}\) by \(\varphi_c(f + \alpha h) = \psi(f) + \alpha c\) for \(f \in U\) and \(\alpha \in \mathbf{R}\) (this is well defined because \(h \notin U\) forces the representation \(f + \alpha h\) of an element of \(U + \mathbf{R}h\) to be unique, and it is clearly linear). Then \(\varphi_c\) is an extension of \(\psi\) with \(\|\varphi_c\| = \|\psi\|\) if and only if \(s(h) \le c \le t(h)\).

Proof of claim. Clearly \(\varphi_c|_U = \psi\), which gives \(\|\varphi_c\| \ge \|\psi\|\); so \(\|\varphi_c\| = \|\psi\|\) if and only if

\begin{equation*} |\psi(f) + \alpha c| \le \|\psi\| \, \|f + \alpha h\| \quad \text{for all } f \in U \text{ and all } \alpha \in \mathbf{R}, \end{equation*}

which is 6.64. As shown in the proof of 6.63, condition 6.64 is equivalent to 6.65, namely

\begin{equation*} |\psi(f) + c| \le \|\psi\| \, \|f + h\| \quad \text{for all } f \in U : \end{equation*}

indeed 6.64 with \(\alpha = 1\) is 6.65, while conversely 6.64 holds trivially for \(\alpha = 0\) (it then reads \(|\psi(f)| \le \|\psi\| \, \|f\|\)) and follows for \(\alpha \ne 0\) by replacing \(f\) with \(\frac{f}{\alpha}\) in 6.65 and multiplying both sides by \(|\alpha|\). Finally, 6.65 states that \(-\|\psi\| \, \|f + h\| \le \psi(f) + c \le \|\psi\| \, \|f + h\|\) for all \(f \in U\), which after subtracting \(\psi(f)\) is

\begin{equation*} -\|\psi\| \, \|f + h\| - \psi(f) \le c \le \|\psi\| \, \|f + h\| - \psi(f) \quad \text{for all } f \in U , \end{equation*}

and taking the supremum over \(f\) on the left and the infimum over \(f\) on the right, this is exactly \(s(h) \le c \le t(h)\). This proves the claim.

Now we prove the two directions.

Suppose first that \(s(h) = t(h)\) for every \(h \in V \setminus U\). The Hahn-Banach Theorem (6.69) gives at least one extension of \(\psi\) to a bounded linear functional on \(V\) whose norm equals \(\|\psi\|\), so we need only prove uniqueness. Suppose \(\varphi_1\) and \(\varphi_2\) are bounded linear functionals on \(V\) with \(\varphi_1|_U = \varphi_2|_U = \psi\) and \(\|\varphi_1\| = \|\varphi_2\| = \|\psi\|\). Let \(h \in V\). If \(h \in U\), then \(\varphi_1(h) = \psi(h) = \varphi_2(h)\). If \(h \in V \setminus U\), set \(c_j = \varphi_j(h)\) for \(j = 1, 2\). For each \(f \in U\) we have

\begin{equation*} |\psi(f) + c_j| = |\varphi_j(f + h)| \le \|\varphi_j\| \, \|f + h\| = \|\psi\| \, \|f + h\| , \end{equation*}

which is 6.65 for \(c_j\); by the computation in the claim this gives \(s(h) \le c_j \le t(h)\). Since \(s(h) = t(h)\), we conclude \(c_1 = s(h) = c_2\), so \(\varphi_1(h) = \varphi_2(h)\). Thus \(\varphi_1 = \varphi_2\), proving that the norm-preserving extension is unique.

Conversely, suppose \(\psi\) has a unique extension to a bounded linear functional on \(V\) of norm \(\|\psi\|\). Suppose, seeking a contradiction, that there exists \(h \in V \setminus U\) with \(s(h) \ne t(h)\); then \(s(h) < t(h)\) by 6.66. Choose \(c_1, c_2 \in [s(h), t(h)]\) with \(c_1 \ne c_2\) (for example \(c_1 = s(h)\) and \(c_2 = t(h)\)). By the claim, for \(j = 1, 2\) the linear functional \(\varphi_{c_j}\) on the subspace \(U + \mathbf{R}h\) of \(V\) extends \(\psi\) and satisfies \(\|\varphi_{c_j}\| = \|\psi\|\). Applying the Hahn-Banach Theorem (6.69) to the bounded linear functional \(\varphi_{c_j}\) on the subspace \(U + \mathbf{R}h\) of \(V\), we obtain a bounded linear functional \(\Phi_j\) on \(V\) with \(\Phi_j|_{U + \mathbf{R}h} = \varphi_{c_j}\) and \(\|\Phi_j\| = \|\varphi_{c_j}\| = \|\psi\|\). Each \(\Phi_j\) restricts to \(\psi\) on \(U\) and has norm \(\|\psi\|\), but

\begin{equation*} \Phi_1(h) = \varphi_{c_1}(h) = c_1 \ne c_2 = \varphi_{c_2}(h) = \Phi_2(h) , \end{equation*}

so \(\Phi_1 \ne \Phi_2\). This contradicts the assumed uniqueness. Hence \(s(h) = t(h)\) for every \(h \in V \setminus U\), completing the proof.

Problem (6D.14)

Show that there exists a linear functional \(\varphi \colon \ell^\infty \to \mathbf{F}\) such that

\begin{equation*} |\varphi(a_1, a_2, \ldots)| \le \|(a_1, a_2, \ldots)\|_\infty \end{equation*}

for all \((a_1, a_2, \ldots) \in \ell^\infty\) and

\begin{equation*} \varphi(a_1, a_2, \ldots) = \lim_{k \to \infty} a_k \end{equation*}

for all \((a_1, a_2, \ldots) \in \ell^\infty\) such that the limit above on the right exists.

Solution

Recall that \(\ell^\infty\) is the normed vector space of bounded sequences in \(\mathbf{F}\), with

\begin{equation*} \|(a_1, a_2, \ldots)\|_\infty = \sup_{k \in \mathbf{Z}^+} |a_k| . \end{equation*}

Let

\begin{equation*} U = \{(a_1, a_2, \ldots) \in \ell^\infty : \lim_{k \to \infty} a_k \text{ exists}\} . \end{equation*}

\(U\) is a subspace of \(\ell^\infty\). Every convergent sequence in \(\mathbf{F}\) is bounded, so a convergent sequence indeed lies in \(\ell^\infty\). The sequence \(0\) converges (to \(0\)), so \(0 \in U\). If \(a = (a_1, a_2, \ldots)\) and \(b = (b_1, b_2, \ldots)\) are in \(U\) and \(\alpha \in \mathbf{F}\), then by the usual limit laws \(\lim_{k \to \infty} (a_k + b_k)\) and \(\lim_{k \to \infty} \alpha a_k\) exist (equalling \(\lim a_k + \lim b_k\) and \(\alpha \lim a_k\)). Hence \(a + b \in U\) and \(\alpha a \in U\).

Define \(\psi \colon U \to \mathbf{F}\) by

\begin{equation*} \psi(a_1, a_2, \ldots) = \lim_{k \to \infty} a_k . \end{equation*}

The limit laws just cited say exactly that \(\psi\) is linear, so \(\psi\) is a linear functional on \(U\).

\(\psi\) is bounded with \(\|\psi\| = 1\). If \(a = (a_1, a_2, \ldots) \in U\) with limit \(L\), then, because the absolute value function on \(\mathbf{F}\) is continuous, \(|L| = \lim_{k \to \infty} |a_k| \le \sup_{k \in \mathbf{Z}^+} |a_k| = \|a\|_\infty\). Hence

\begin{equation*} |\psi(a)| \le \|a\|_\infty \quad \text{for all } a \in U, \end{equation*}

so \(\psi\) is a bounded linear functional with \(\|\psi\| \le 1\). Moreover the constant sequence \(u = (1, 1, 1, \ldots)\) belongs to \(U\) with \(\|u\|_\infty = 1\) and \(\psi(u) = 1\), so \(\|\psi\| \ge 1\). Thus \(\|\psi\| = 1\).

By the Hahn-Banach Theorem (6.69), applied to the subspace \(U\) of the normed vector space \(\ell^\infty\) and the bounded linear functional \(\psi\) on \(U\), there exists a bounded linear functional \(\varphi\) on \(\ell^\infty\) such that

\begin{equation*} \varphi|_U = \psi \quad \text{and} \quad \|\varphi\| = \|\psi\| = 1 . \end{equation*}

This \(\varphi\) has the two required properties. First, for every \((a_1, a_2, \ldots) \in \ell^\infty\), the definition of the norm of a bounded linear functional gives

\begin{equation*} |\varphi(a_1, a_2, \ldots)| \le \|\varphi\| \, \|(a_1, a_2, \ldots)\|_\infty = \|(a_1, a_2, \ldots)\|_\infty . \end{equation*}

Second, if \((a_1, a_2, \ldots) \in \ell^\infty\) is such that \(\lim_{k \to \infty} a_k\) exists, then \((a_1, a_2, \ldots) \in U\) and hence

\begin{equation*} \varphi(a_1, a_2, \ldots) = \psi(a_1, a_2, \ldots) = \lim_{k \to \infty} a_k . \end{equation*}

Thus a linear functional \(\varphi\) with the required properties exists. Such a norm-one extension of the limit functional is the starting point for the construction of a Banach limit (which is additionally required to be invariant under the shift operator). Note that the existence of \(\varphi\) came from the Hahn-Banach Theorem, whose proof uses Zorn’s Lemma; the argument gives no formula for \(\varphi\) on a sequence that does not converge.

Problem (6D.15)

Suppose \(B\) is an open ball in a normed vector space \(V\) such that \(0 \notin B\). Prove that there exists \(\varphi \in V^{\prime}\) such that

\begin{equation*} \operatorname{Re} \varphi(f) > 0 \end{equation*}

for all \(f \in B\).

Solution

Write the open ball as

\begin{equation*} B = \{ f \in V : \| f - g \| < r \}, \end{equation*}

where \(g \in V\) is its center and \(r > 0\) is its radius.

Because \(0 \notin B\), we have \(\| 0 - g \| \geq r\), which means \(\| g \| \geq r > 0\). In particular \(g \neq 0\).

Because \(g \in V \setminus \{0\}\), the Hahn-Banach Theorem in the form 6.72 provides \(\varphi \in V^{\prime}\) with

\begin{equation*} \| \varphi \| = 1 \qquad \text{and} \qquad \varphi(g) = \| g \|. \end{equation*}

Now suppose \(f \in B\). Then

\begin{equation*} \begin{aligned} \operatorname{Re} \varphi(f) &= \operatorname{Re} \varphi(g) + \operatorname{Re} \varphi(f - g) \\ &= \| g \| + \operatorname{Re} \varphi(f - g) \\ &\geq \| g \| - | \varphi(f - g) | \\ &\geq \| g \| - \| \varphi \| \, \| f - g \| \\ &= \| g \| - \| f - g \| \\ &> \| g \| - r \\ &\geq 0, \end{aligned} \end{equation*}

where the second line uses \(\varphi(g) = \| g \| \in \mathbb{R}\), the third line uses \(\operatorname{Re} z \geq -|z|\), the fourth line uses the definition of \(\| \varphi \|\), the sixth line uses \(\| f - g \| < r\) (which is the statement that \(f \in B\)), and the last line uses \(\| g \| \geq r\).

Thus \(\operatorname{Re} \varphi(f) > 0\) for every \(f \in B\), as desired.

Note that if \(\mathbb{F} = \mathbb{R}\) then \(\operatorname{Re} \varphi = \varphi\), so the conclusion reads \(\varphi(f) > 0\) for all \(f \in B\). Geometrically, the result says that an open ball missing the origin lies strictly on one side of a closed hyperplane through the origin, namely the hyperplane \(\{ f \in V : \operatorname{Re} \varphi(f) = 0 \}\).

Problem (6D.16)

Show that the dual space of each infinite-dimensional normed vector space is infinite-dimensional.

Solution

Suppose \(V\) is an infinite-dimensional normed vector space. We show that for every \(n \in \mathbb{Z}^+\) the dual space \(V^{\prime}\) contains a linearly independent list of length \(n\); because a finite-dimensional vector space of dimension \(m\) contains no linearly independent list of length \(m + 1\), this implies that \(V^{\prime}\) is infinite-dimensional.

Step 1: \(V\) contains a linearly independent list of length \(n\) for each \(n \in \mathbb{Z}^+\).

We construct such a list inductively. Because \(V\) is infinite-dimensional, \(V \neq \{0\}\), so there exists \(f_1 \in V \setminus \{0\}\); the list \(f_1\) is linearly independent. Suppose \(k \in \mathbb{Z}^+\) and linearly independent \(f_1, \ldots, f_k \in V\) have been found. Then \(\operatorname{span}\{f_1, \ldots, f_k\} \neq V\), because otherwise \(V\) would be finite-dimensional (see 6.54). Hence there exists

\begin{equation*} f_{k+1} \in V \setminus \operatorname{span}\{f_1, \ldots, f_k\}. \end{equation*}

If \(\alpha_1 f_1 + \cdots + \alpha_{k+1} f_{k+1} = 0\) with not all \(\alpha_j\) equal to \(0\), then \(\alpha_{k+1} = 0\) (otherwise \(f_{k+1}\) would be in the span of \(f_1, \ldots, f_k\)), and then the linear independence of \(f_1, \ldots, f_k\) forces \(\alpha_1 = \cdots = \alpha_k = 0\). Thus \(f_1, \ldots, f_{k+1}\) is linearly independent, completing the induction.

Step 2: producing \(n\) linearly independent bounded linear functionals.

Fix \(n \in \mathbb{Z}^+\) and let \(f_1, \ldots, f_n \in V\) be linearly independent, as provided by Step 1. Let

\begin{equation*} W = \operatorname{span}\{f_1, \ldots, f_n\}, \end{equation*}

which is a finite-dimensional subspace of \(V\); regard \(W\) as a normed vector space with the norm it inherits from \(V\). Because \(f_1, \ldots, f_n\) is linearly independent and spans \(W\), every element of \(W\) has a unique representation \(\alpha_1 f_1 + \cdots + \alpha_n f_n\) with \(\alpha_1, \ldots, \alpha_n \in \mathbb{F}\). Hence for each \(j \in \{1, \ldots, n\}\) we get a well-defined linear functional \(\psi_j : W \to \mathbb{F}\) (the \(j\)-th coordinate functional) defined by

\begin{equation*} \psi_j(\alpha_1 f_1 + \cdots + \alpha_n f_n) = \alpha_j . \end{equation*}

Because \(W\) is finite-dimensional, every linear map from \(W\) to a normed vector space is continuous (Exercise 7 in Section 6D); thus each \(\psi_j\) is a bounded linear functional on \(W\).

By the Hahn-Banach Theorem (6.69), each \(\psi_j\) extends to a bounded linear functional \(\varphi_j \in V^{\prime}\) (with \(\| \varphi_j \| = \| \psi_j \|\), although we do not need the norm equality here). By construction,

\begin{equation*} \varphi_j(f_k) = \psi_j(f_k) = \begin{cases} 1 & \text{if } j = k, \\ 0 & \text{if } j \neq k, \end{cases} \end{equation*}

for all \(j, k \in \{1, \ldots, n\}\).

Step 3: \(\varphi_1, \ldots, \varphi_n\) is linearly independent in \(V^{\prime}\).

Suppose \(\alpha_1, \ldots, \alpha_n \in \mathbb{F}\) and \(\alpha_1 \varphi_1 + \cdots + \alpha_n \varphi_n = 0\). Evaluating this functional at \(f_k\) and using \(\varphi_j(f_k) = 0\) for \(j \neq k\) and \(\varphi_k(f_k) = 1\) gives \(\alpha_k = 0\), for each \(k \in \{1, \ldots, n\}\). Thus \(\varphi_1, \ldots, \varphi_n\) is a linearly independent list in \(V^{\prime}\).

Because \(n \in \mathbb{Z}^+\) was arbitrary, \(V^{\prime}\) contains linearly independent lists of every finite length. If \(V^{\prime}\) were finite-dimensional, say spanned by a list of length \(m\), then every list in \(V^{\prime}\) of length \(m + 1\) would be linearly dependent, contradicting what we just proved. Hence \(V^{\prime}\) is infinite-dimensional.

Problem (6D.17)

Suppose \(V\) is a separable normed vector space. Explain how the Hahn-Banach Theorem (6.69) for \(V\) can be proved without using any results (such as Zorn’s Lemma) that depend on the Axiom of Choice.

Solution

Zorn’s Lemma enters the proof of 6.69 only to produce a maximal extension. When \(V\) is separable we can instead extend one basis element at a time along a fixed enumeration of a countable dense set, then extend by continuity. The only thing that needs care is that every step must be given by an explicit rule, so that the sequence of extensions is produced by ordinary recursion rather than by infinitely many arbitrary choices.

Throughout, let \(g_1, g_2, \ldots\) be a fixed enumeration of a countable dense subset \(G\) of \(V\), indexed by all of \(\mathbb{Z}^+\) (if \(G\) is finite and nonempty, repeat its elements; if \(G = \varnothing\) then \(V = \{0\}\) and the result is trivial). Such an enumeration exists by the separability hypothesis, and fixing one is a single existential instantiation, not an application of the Axiom of Choice. Let \(U\) be a subspace of \(V\) and let \(\psi : U \to \mathbb{F}\) be a bounded linear functional.

Case \(\mathbb{F} = \mathbb{R}\).

First note that the Extension Lemma (6.63) uses no choice principle once we specify \(c\) explicitly. Indeed, in the setting of 6.63 (with \(\psi\) replaced by a bounded linear functional \(\rho\) on a subspace \(W\) of \(V\), and \(h \in V \setminus W\)), define

\begin{equation*} c = \sup_{f \in W} \bigl( -\| \rho \| \, \| f + h \| - \rho(f) \bigr). \end{equation*}

This supremum is a well-defined real number: taking \(f = 0\) shows it is at least \(-\| \rho \| \, \| h \|\), and inequality 6.66 shows it is at most \(\inf_{g \in W} ( \| \rho \| \, \| g + h \| - \rho(g) ) \leq \| \rho \| \, \| h \|\). With this specific \(c\), the computation in the proof of 6.63 shows that

\begin{equation*} \rho_h(f + \alpha h) = \rho(f) + \alpha c \qquad (f \in W, \ \alpha \in \mathbb{R}) \end{equation*}

defines an extension of \(\rho\) to \(W + \mathbb{R}h\) with \(\| \rho_h \| = \| \rho \|\). Thus we have an explicit function assigning to each pair \((W, \rho)\) and each \(h \in V \setminus W\) a norm-preserving extension of \(\rho\) to \(W + \mathbb{R}h\).

Now define recursively \(U_0 = U\) and \(\varphi_0 = \psi\), and, given a subspace \(U_n\) of \(V\) and a bounded linear functional \(\varphi_n\) on \(U_n\) with \(\| \varphi_n \| = \| \psi \|\):

  • if \(g_k \in U_n\) for every \(k \in \mathbb{Z}^+\), set \(U_{n+1} = U_n\) and \(\varphi_{n+1} = \varphi_n\);
  • otherwise let \(m_n\) be the smallest element of \(\mathbb{Z}^+\) with \(g_{m_n} \notin U_n\), set \(U_{n+1} = U_n + \mathbb{R} g_{m_n}\), and let \(\varphi_{n+1}\) be the extension of \(\varphi_n\) to \(U_{n+1}\) given by the explicit Extension Lemma above (with \(h = g_{m_n}\)).

Every ingredient here is uniquely determined by \((U_n, \varphi_n)\) and the fixed enumeration: the index \(m_n\) is a least element of a nonempty subset of \(\mathbb{Z}^+\), and the constant \(c\) is a supremum. Hence the sequence \((U_n, \varphi_n)_{n \geq 0}\) is produced by the recursion theorem, with no choices made.

Each \(\varphi_{n+1}\) extends \(\varphi_n\) and \(\| \varphi_n \| = \| \psi \|\) for every \(n\) (the norm is preserved at each step, and \(\| \varphi_n \| \geq \| \psi \|\) because \(\varphi_n\) extends \(\psi\)).

Next, \(\{ g_1, g_2, \ldots \} \subseteq \bigcup_{n=0}^\infty U_n\). If the first bullet point above ever applies, this is immediate. Otherwise \(m_n\) is defined for every \(n\), and since \(g_{m_n} \in U_{n+1}\) while \(U_n \subseteq U_{n+1}\), the least missing index strictly increases: \(m_{n+1} > m_n\). Hence \(m_n \geq m_0 + n \to \infty\), so for each \(k \in \mathbb{Z}^+\) there exists \(n\) with \(m_n > k\), which forces \(g_k \in U_n\).

Let \(W = \bigcup_{n=0}^\infty U_n\). Because \(U_0 \subseteq U_1 \subseteq \cdots\), the set \(W\) is a subspace of \(V\), and it contains the dense set \(\{ g_1, g_2, \ldots \}\); hence \(W\) is dense in \(V\). Define \(\varphi : W \to \mathbb{R}\) by \(\varphi(f) = \varphi_n(f)\) for any \(n\) with \(f \in U_n\); this is well defined because the \(\varphi_n\) are successive extensions of one another. Then \(\varphi\) is linear (any two elements of \(W\) lie in a common \(U_n\)), \(\varphi\) extends \(\psi\), and \(| \varphi(f) | \leq \| \psi \| \, \| f \|\) for all \(f \in W\).

Finally we extend \(\varphi\) from the dense subspace \(W\) to all of \(V\), again without choices. For \(f \in V\) and \(n \in \mathbb{Z}^+\), let \(k_n(f)\) be the smallest \(k \in \mathbb{Z}^+\) with \(\| g_k - f \| < \frac{1}{n}\); such \(k\) exists by density, and \(g_{k_n(f)} \in W\). For \(m, n \in \mathbb{Z}^+\),

\begin{equation*} \bigl| \varphi(g_{k_n(f)}) - \varphi(g_{k_m(f)}) \bigr| \leq \| \psi \| \, \| g_{k_n(f)} - g_{k_m(f)} \| \leq \| \psi \| \Bigl( \frac{1}{n} + \frac{1}{m} \Bigr), \end{equation*}

so \(\bigl( \varphi(g_{k_n(f)}) \bigr)_{n \in \mathbb{Z}^+}\) is a Cauchy sequence of real numbers. Define

\begin{equation*} \tilde{\varphi}(f) = \lim_{n \to \infty} \varphi(g_{k_n(f)}), \end{equation*}

which exists by the completeness of \(\mathbb{R}\) (a fact requiring no choice, since the limit can be written as \(\lim_n \varphi(g_{k_n(f)}) = \sup_n \inf_{j \geq n} \varphi(g_{k_j(f)})\)).

If \(h_1, h_2, \ldots\) is any sequence in \(W\) with \(h_j \to f\), then

\begin{equation*} | \varphi(h_j) - \varphi(g_{k_n(f)}) | \leq \| \psi \| \bigl( \| h_j - f \| + \tfrac{1}{n} \bigr), \end{equation*}

so \(\varphi(h_j) \to \tilde{\varphi}(f)\). Applying this to \(f, f^{\prime} \in V\), \(\alpha \in \mathbb{R}\) and sequences \(h_j \to f\), \(h_j^{\prime} \to f^{\prime}\) in \(W\) (for instance \(h_j = g_{k_j(f)}\), \(h_j^{\prime} = g_{k_j(f^{\prime})}\), so no choice is needed), and using \(h_j + \alpha h_j^{\prime} \to f + \alpha f^{\prime}\), we get

\begin{equation*} \tilde{\varphi}(f + \alpha f^{\prime}) = \lim_{j \to \infty} \varphi(h_j + \alpha h_j^{\prime}) = \lim_{j \to \infty} \bigl( \varphi(h_j) + \alpha \varphi(h_j^{\prime}) \bigr) = \tilde{\varphi}(f) + \alpha \tilde{\varphi}(f^{\prime}). \end{equation*}

Thus \(\tilde{\varphi}\) is a linear functional on \(V\). Taking \(f \in W\) and \(h_j = f\) for all \(j\) shows \(\tilde{\varphi}|_W = \varphi\), so \(\tilde{\varphi}\) extends \(\psi\). Also \(| \varphi(g_{k_n(f)}) | \leq \| \psi \| \, \| g_{k_n(f)} \|\) and \(g_{k_n(f)} \to f\) give \(| \tilde{\varphi}(f) | \leq \| \psi \| \, \| f \|\) for all \(f \in V\). Hence \(\| \tilde{\varphi} \| \leq \| \psi \|\), and \(\| \tilde{\varphi} \| \geq \| \psi \|\) because \(\tilde{\varphi}\) extends \(\psi\). Therefore \(\| \tilde{\varphi} \| = \| \psi \|\), which is the Hahn-Banach Theorem for \(V\) in the real case.

Case \(\mathbb{F} = \mathbb{C}\).

The reduction to the real case given in the proof of 6.69 uses no choice principle at all: one sets \(\psi_1 = \operatorname{Re} \psi\) (an \(\mathbb{R}\)-linear functional on \(U\) with \(\| \psi_1 \| \leq \| \psi \|\)), regards \(V\) as a real normed vector space, applies the real case above to obtain an \(\mathbb{R}\)-linear \(\varphi_1 : V \to \mathbb{R}\) extending \(\psi_1\) with \(\| \varphi_1 \| = \| \psi_1 \|\), and then defines

\begin{equation*} \varphi(f) = \varphi_1(f) - i \varphi_1(if) \end{equation*}

for \(f \in V\). Note that \(V\) regarded as a real normed vector space has exactly the same norm and hence exactly the same countable dense subset \(\{ g_1, g_2, \ldots \}\), so it is separable and the real case does apply. The verifications in the proof of 6.69 that \(\varphi\) is \(\mathbb{C}\)-linear, extends \(\psi\), and satisfies \(\| \varphi \| \leq \| \psi \|\) are pure computations.

In summary: separability lets us replace the maximal element supplied by Zorn’s Lemma by an explicitly defined increasing sequence of one-dimensional extensions along a fixed enumeration of a countable dense set, followed by an explicitly defined extension by continuity. Every step is given by a formula (a least index, a supremum, a limit), so the whole argument is carried out in ZF.

Problem (6D.18)

Suppose \(V\) is a normed vector space such that the dual space \(V^{\prime}\) is a separable Banach space. Prove that \(V\) is separable.

Solution

By 6.47 the dual space \(V^{\prime}\) is always a Banach space, so the hypothesis with content is that \(V^{\prime}\) is separable.

If \(V = \{0\}\) then \(V\) is separable and there is nothing to prove, so assume \(V \neq \{0\}\).

Let \(\{ \varphi_1, \varphi_2, \ldots \}\) be a countable dense subset of \(V^{\prime}\). Discarding those \(\varphi_n\) that equal \(0\) changes nothing (the remaining set is still dense in \(V^{\prime}\), because \(0\) is a limit of nonzero elements of \(V^{\prime}\): by 6.72 there exists \(\varphi \in V^{\prime}\) with \(\| \varphi \| = 1\), and then \(\varphi / n \to 0\); alternatively, simply keep all \(\varphi_n\) and note that the argument below only ever uses those with \(\| \varphi_n \| \geq \frac{3}{4}\)). So assume \(\varphi_n \neq 0\) for every \(n \in \mathbb{Z}^+\).

For each \(n \in \mathbb{Z}^+\), the definition of \(\| \varphi_n \|\) as a supremum gives some \(f_n \in V\) with

\begin{equation*} \| f_n \| \leq 1 \qquad \text{and} \qquad | \varphi_n(f_n) | > \tfrac{1}{2} \| \varphi_n \| ; \end{equation*}

fix such an \(f_n\) for each \(n\).

Let

\begin{equation*} W = \operatorname{span}\{ f_n : n \in \mathbb{Z}^+ \}, \end{equation*}

the set of all finite linear combinations of the \(f_n\), and let

\begin{equation*} D = \Bigl\{ \sum_{n \in \Omega} \alpha_n f_n : \Omega \subseteq \mathbb{Z}^+ \text{ finite and } \alpha_n \in \mathbb{Q} \text{ for all } n \in \Omega \Bigr\} \end{equation*}

if \(\mathbb{F} = \mathbb{R}\), and the same with \(\mathbb{Q}\) replaced by \(\mathbb{Q} + i\mathbb{Q}\) if \(\mathbb{F} = \mathbb{C}\). Then \(D\) is a countable set (a countable union, over the countably many finite subsets \(\Omega\) of \(\mathbb{Z}^+\), of countable sets).

We claim \(\overline{D} = V\); this proves that \(V\) is separable.

First, \(W \subseteq \overline{D}\). Indeed, if \(\Omega \subseteq \mathbb{Z}^+\) is finite and \(\alpha_n \in \mathbb{F}\) for \(n \in \Omega\), then choosing \(\beta_n\) in \(\mathbb{Q}\) (or \(\mathbb{Q} + i\mathbb{Q}\)) with \(| \alpha_n - \beta_n |\) small gives

\begin{equation*} \Bigl\| \sum_{n \in \Omega} \alpha_n f_n - \sum_{n \in \Omega} \beta_n f_n \Bigr\| \leq \sum_{n \in \Omega} | \alpha_n - \beta_n | \, \| f_n \| \leq \sum_{n \in \Omega} | \alpha_n - \beta_n |, \end{equation*}

which can be made arbitrarily small. Hence \(\overline{D} \supseteq \overline{W}\).

So it suffices to prove that \(\overline{W} = V\). Suppose not. Then there exists \(h \in V\) with \(h \notin \overline{W}\). Because \(W\) is a subspace of \(V\), the result 6.73 implies that there exists \(\varphi \in V^{\prime}\) with

\begin{equation*} \varphi|_W = 0 \qquad \text{and} \qquad \varphi(h) \neq 0 . \end{equation*}

In particular \(\varphi \neq 0\), so after replacing \(\varphi\) by \(\varphi / \| \varphi \|\) we may assume \(\| \varphi \| = 1\).

Because \(\{ \varphi_1, \varphi_2, \ldots \}\) is dense in \(V^{\prime}\), there exists \(n \in \mathbb{Z}^+\) with

\begin{equation*} \| \varphi - \varphi_n \| < \tfrac{1}{4} . \end{equation*}

Then

\begin{equation*} \| \varphi_n \| \geq \| \varphi \| - \| \varphi - \varphi_n \| > 1 - \tfrac{1}{4} = \tfrac{3}{4} . \end{equation*}

On the other hand, \(f_n \in W\) and \(\varphi|_W = 0\) give \(\varphi(f_n) = 0\), so

\begin{equation*} \tfrac{1}{2} \| \varphi_n \| < | \varphi_n(f_n) | = | \varphi_n(f_n) - \varphi(f_n) | \leq \| \varphi_n - \varphi \| \, \| f_n \| < \tfrac{1}{4} \cdot 1 = \tfrac{1}{4}, \end{equation*}

which gives \(\| \varphi_n \| < \frac{1}{2}\). This contradicts \(\| \varphi_n \| > \frac{3}{4}\).

Hence no such \(h\) exists, so \(\overline{W} = V\) and therefore \(\overline{D} = V\). Because \(D\) is countable, \(V\) is separable.

Remark: the converse fails. The Banach space \(C([0,1])\) is separable (by the Weierstrass approximation theorem the polynomials are dense in it, and the polynomials whose coefficients lie in \(\mathbb{Q}\), or in \(\mathbb{Q} + i\mathbb{Q}\) when \(\mathbb{F} = \mathbb{C}\), form a countable dense subset), but its dual is not separable, by the next exercise.

Problem (6D.19)

Prove that the dual of the Banach space \(C([0,1])\) is not separable; here the norm on \(C([0,1])\) is defined by \(\| f \| = \sup_{[0,1]} | f |\).

Solution

For each \(t \in [0,1]\) define \(\varphi_t : C([0,1]) \to \mathbb{F}\) by

\begin{equation*} \varphi_t(f) = f(t) . \end{equation*}

Each \(\varphi_t\) is clearly linear, and

\begin{equation*} | \varphi_t(f) | = | f(t) | \leq \sup_{[0,1]} | f | = \| f \|, \end{equation*}

so \(\varphi_t\) is a bounded linear functional with \(\| \varphi_t \| \leq 1\). Taking \(f\) to be the constant function \(1\) (which has norm \(1\)) gives \(\varphi_t(f) = 1\), so in fact \(\| \varphi_t \| = 1\). Thus \(\varphi_t \in C([0,1])^{\prime}\) for every \(t \in [0,1]\).

Claim: if \(s, t \in [0,1]\) and \(s \neq t\), then \(\| \varphi_s - \varphi_t \| = 2\).

To see this, suppose \(s \neq t\) and set \(\delta = | s - t | > 0\). Define \(f : [0,1] \to \mathbb{R}\) by

\begin{equation*} f(x) = \max \Bigl\{ -1, \ 1 - \frac{2 | x - s |}{\delta} \Bigr\} . \end{equation*}

Then \(f\) is continuous (a maximum of two continuous functions), \(-1 \leq f \leq 1\) so \(\| f \| \leq 1\), and \(f(s) = 1\) while \(f(t) = \max\{-1, 1 - 2\} = -1\). Hence

\begin{equation*} \| \varphi_s - \varphi_t \| \geq | (\varphi_s - \varphi_t)(f) | = | f(s) - f(t) | = 2 . \end{equation*}

The reverse inequality \(\| \varphi_s - \varphi_t \| \leq \| \varphi_s \| + \| \varphi_t \| = 2\) holds by the triangle inequality in \(C([0,1])^{\prime}\). This proves the claim.

Now suppose, seeking a contradiction, that \(C([0,1])^{\prime}\) is separable, and let \(\{ \mu_1, \mu_2, \ldots \}\) be a countable dense subset of \(C([0,1])^{\prime}\).

For each \(t \in [0,1]\) there exists \(j \in \mathbb{Z}^+\) with \(\| \mu_j - \varphi_t \| < 1\) (by density). Define \(\kappa : [0,1] \to \mathbb{Z}^+\) by letting \(\kappa(t)\) be the smallest such \(j\); this is a well-defined function, requiring no choice, because every nonempty subset of \(\mathbb{Z}^+\) has a smallest element.

The function \(\kappa\) is one-to-one. Indeed, suppose \(s, t \in [0,1]\) and \(\kappa(s) = \kappa(t) = j\). Then

\begin{equation*} \| \varphi_s - \varphi_t \| \leq \| \varphi_s - \mu_j \| + \| \mu_j - \varphi_t \| < 1 + 1 = 2 . \end{equation*}

By the claim above, this forces \(s = t\).

Thus \(\kappa\) is a one-to-one function from the uncountable set \([0,1]\) into the countable set \(\mathbb{Z}^+\), which is impossible.

Therefore \(C([0,1])^{\prime}\) has no countable dense subset; that is, the dual of \(C([0,1])\) is not separable.

Remark: the same argument shows more generally that if a normed vector space contains an uncountable subset whose elements are pairwise at distance at least some fixed \(\varepsilon > 0\), then that space is not separable. Note also that \(C([0,1])\) itself is separable, so this exercise shows that the converse of Exercise 18 is false: a separable normed vector space can have a nonseparable dual.

Problem (6D.20)

Define \(\Phi : V \to V^{\prime\prime}\) by

\begin{equation*} (\Phi f)(\varphi) = \varphi(f) \end{equation*}

for \(f \in V\) and \(\varphi \in V^{\prime}\). Show that \(\| \Phi f \| = \| f \|\) for every \(f \in V\).

[The map \(\Phi\) defined above is called the canonical isometry of \(V\) into \(V^{\prime\prime}\).]

Solution

Here \(V\) is a normed vector space, \(V^{\prime} = \mathcal{B}(V, \mathbb{F})\) is its dual space, and \(V^{\prime\prime} = (V^{\prime})^{\prime}\) is its double dual space, with the norm \(V^{\prime\prime}\) receives as the dual space of \(V^{\prime}\).

First we check that \(\Phi\) really maps into \(V^{\prime\prime}\), that is, that \(\Phi f\) is a bounded linear functional on \(V^{\prime}\) for each \(f \in V\).

Linearity: if \(\varphi, \psi \in V^{\prime}\) and \(\alpha \in \mathbb{F}\), then by the definition of addition and scalar multiplication in \(V^{\prime}\),

\begin{equation*} (\Phi f)(\varphi + \alpha \psi) = (\varphi + \alpha \psi)(f) = \varphi(f) + \alpha \psi(f) = (\Phi f)(\varphi) + \alpha (\Phi f)(\psi) . \end{equation*}

Thus \(\Phi f : V^{\prime} \to \mathbb{F}\) is linear.

Boundedness, and the inequality \(\| \Phi f \| \leq \| f \|\): for every \(\varphi \in V^{\prime}\),

\begin{equation*} | (\Phi f)(\varphi) | = | \varphi(f) | \leq \| \varphi \| \, \| f \| , \end{equation*}

where the inequality is the definition of \(\| \varphi \|\). Hence \(\Phi f\) is a bounded linear functional on \(V^{\prime}\) with

\begin{equation*} \| \Phi f \| \leq \| f \| . \end{equation*}

In particular \(\Phi f \in V^{\prime\prime}\), so \(\Phi\) is indeed a map from \(V\) to \(V^{\prime\prime}\).

Now we prove the reverse inequality \(\| \Phi f \| \geq \| f \|\).

If \(f = 0\) then \(\Phi f = 0\) and both norms equal \(0\), so the desired equality holds. Hence assume \(f \in V \setminus \{0\}\). By 6.72 (a consequence of the Hahn-Banach Theorem 6.69), there exists \(\varphi \in V^{\prime}\) with

\begin{equation*} \| \varphi \| = 1 \qquad \text{and} \qquad \varphi(f) = \| f \| . \end{equation*}

Then

\begin{equation*} \| \Phi f \| \geq | (\Phi f)(\varphi) | = | \varphi(f) | = \| f \| , \end{equation*}

where the first inequality holds because \(\| \Phi f \|\) is the supremum of \(| (\Phi f)(\psi) |\) over all \(\psi \in V^{\prime}\) with \(\| \psi \| \leq 1\), and \(\| \varphi \| = 1\).

Combining the two inequalities gives \(\| \Phi f \| = \| f \|\) for every \(f \in V\), as desired.

Finally, note that \(\Phi\) is linear: for \(f, g \in V\), \(\alpha \in \mathbb{F}\), and every \(\varphi \in V^{\prime}\),

\begin{equation*} \bigl( \Phi(f + \alpha g) \bigr)(\varphi) = \varphi(f + \alpha g) = \varphi(f) + \alpha \varphi(g) = (\Phi f)(\varphi) + \alpha (\Phi g)(\varphi) = \bigl( \Phi f + \alpha \, \Phi g \bigr)(\varphi), \end{equation*}

so \(\Phi(f + \alpha g) = \Phi f + \alpha \, \Phi g\). Thus \(\Phi\) is a linear map that preserves norms; in particular \(\Phi\) is one-to-one (if \(\Phi f = 0\) then \(\| f \| = \| \Phi f \| = 0\), so \(f = 0\)), which justifies calling \(\Phi\) the canonical isometry of \(V\) into \(V^{\prime\prime}\).

Problem (6D.21)

Suppose \(V\) is an infinite-dimensional normed vector space. Show that there is a convex subset \(U\) of \(V\) such that \(\overline{U} = V\) and such that the complement \(V \setminus U\) is also a convex subset of \(V\) with \(\overline{V \setminus U} = V\).

[See 8.25 for the definition of a convex set. This exercise should stretch your geometric intuition because this behavior cannot happen in finite dimensions.]

Solution

Recall from 8.25 that a subset \(U\) of a vector space is convex if \((1-t) f + t g \in U\) for all \(t \in [0,1]\) and all \(f, g \in U\). Convexity involves only real scalars, so it is convenient to work with \(V\) regarded as a real vector space.

Step 1: \(V\) is infinite-dimensional as a real normed vector space.

If \(\mathbb{F} = \mathbb{R}\) there is nothing to prove. If \(\mathbb{F} = \mathbb{C}\), let \(V_{\mathbb{R}}\) denote \(V\) with scalar multiplication restricted to \(\mathbb{R}\), equipped with the same norm; this is a real normed vector space. If \(V_{\mathbb{R}}\) were finite-dimensional, say \(V_{\mathbb{R}} = \operatorname{span}_{\mathbb{R}}\{f_1, \ldots, f_n\}\), then also \(V = \operatorname{span}_{\mathbb{C}}\{f_1, \ldots, f_n\}\), contradicting the hypothesis that \(V\) is infinite-dimensional. Hence \(V_{\mathbb{R}}\) is infinite-dimensional. From now on we write \(V\) for \(V_{\mathbb{R}}\) and consider only real scalars.

Step 2: choose an unbounded \(\mathbb{R}\)-linear functional.

By 6.62 applied to the infinite-dimensional real normed vector space \(V\), there exists a discontinuous linear functional \(\rho : V \to \mathbb{R}\). By 6.48 (or directly by 6.52), discontinuous means unbounded: for every \(c \in [0, \infty)\) there exists \(f \in V\) with \(| \rho(f) | > c \| f \|\). In particular \(\rho\) is not identically \(0\).

Step 3: \(\operatorname{null} \rho\) is dense in \(V\).

Because \(\rho\) is unbounded, for each \(j \in \mathbb{Z}^+\) there exists \(h_j \in V\) with \(\| h_j \| \leq 1\) and \(| \rho(h_j) | > j\) (start from \(f_j\) with \(| \rho(f_j) | > j \| f_j \|\), note \(f_j \neq 0\), and set \(h_j = f_j / \| f_j \|\)). In particular \(\rho(h_j) \neq 0\).

Now suppose \(g \in V\). For each \(j \in \mathbb{Z}^+\) define

\begin{equation*} k_j = g - \frac{\rho(g)}{\rho(h_j)} h_j . \end{equation*}

Then \(\rho(k_j) = \rho(g) - \rho(g) = 0\), so \(k_j \in \operatorname{null} \rho\), and

\begin{equation*} \| k_j - g \| = \frac{| \rho(g) |}{| \rho(h_j) |} \| h_j \| \leq \frac{| \rho(g) |}{j} \longrightarrow 0 \end{equation*}

as \(j \to \infty\). Hence \(g \in \overline{\operatorname{null} \rho}\). Because \(g \in V\) was arbitrary, \(\operatorname{null} \rho\) is dense in \(V\). (This is exactly the implication (d) \(\Rightarrow\) (a) of 6.52, read in contrapositive form.)

Step 4: the set \(U\) and its complement.

Define

\begin{equation*} U = \{ f \in V : \rho(f) > 0 \}, \qquad \text{so} \qquad V \setminus U = \{ f \in V : \rho(f) \leq 0 \} . \end{equation*}

Both are convex. Indeed, suppose \(f, g \in U\) and \(t \in [0,1]\). Then

\begin{equation*} \rho\bigl( (1-t) f + t g \bigr) = (1-t) \rho(f) + t \rho(g) > 0, \end{equation*}

because \((1-t) \rho(f) \geq 0\) and \(t \rho(g) \geq 0\) while at least one of \(1 - t\) and \(t\) is nonzero (and \(\rho(f), \rho(g) > 0\)), so at least one of the two terms is strictly positive. Hence \((1-t) f + t g \in U\), so \(U\) is convex. Similarly, if \(f, g \in V \setminus U\) and \(t \in [0,1]\), then

\begin{equation*} \rho\bigl( (1-t) f + t g \bigr) = (1-t) \rho(f) + t \rho(g) \leq 0, \end{equation*}

because both terms are less than or equal to \(0\); hence \((1-t) f + t g \in V \setminus U\), so \(V \setminus U\) is convex.

Step 5: both sets are dense in \(V\).

Because \(\operatorname{null} \rho \subseteq V \setminus U\) and \(\operatorname{null} \rho\) is dense in \(V\) by Step 3, we get \(\overline{V \setminus U} = V\).

For \(U\), first fix \(h \in V\) with \(\rho(h) = 1\): because \(\rho\) is not identically \(0\), there exists \(h_0 \in V\) with \(\rho(h_0) \neq 0\), and we may take \(h = h_0 / \rho(h_0)\). Note \(h \neq 0\), so \(\| h \| > 0\).

Now suppose \(g \in V\) and \(\varepsilon > 0\). By Step 3 there exists \(k \in \operatorname{null} \rho\) with \(\| k - g \| < \varepsilon / 2\). Let

\begin{equation*} \delta = \frac{\varepsilon}{2 \| h \| + 2} > 0 \qquad \text{and} \qquad u = k + \delta h . \end{equation*}

Then \(\rho(u) = \rho(k) + \delta \rho(h) = \delta > 0\), so \(u \in U\), and

\begin{equation*} \| u - g \| \leq \| k - g \| + \delta \| h \| < \frac{\varepsilon}{2} + \frac{\varepsilon \| h \|}{2 \| h \| + 2} < \frac{\varepsilon}{2} + \frac{\varepsilon}{2} = \varepsilon . \end{equation*}

Because \(g \in V\) and \(\varepsilon > 0\) were arbitrary, \(\overline{U} = V\).

Thus \(U\) is a convex subset of \(V\) with \(\overline{U} = V\) whose complement \(V \setminus U\) is also convex with \(\overline{V \setminus U} = V\), as desired. (Note that \(U\) and \(V \setminus U\) are convex as subsets of \(V\) regardless of whether \(\mathbb{F}\) is \(\mathbb{R}\) or \(\mathbb{C}\), since convexity refers only to real coefficients, and density is with respect to the norm of \(V\), which is unchanged by restricting scalars.)

Remark on why this cannot happen in finite dimensions.

Suppose \(V\) is finite-dimensional, say of real dimension \(n\), and \(U \subseteq V\) is convex with \(\overline{U} = V\). We claim \(U\) has nonempty interior. Identify \(V\) with \(\mathbb{R}^n\) (all norms on \(\mathbb{R}^n\) give the same open sets by Exercise 6 in Section 6D). The points \(0, e_1, \ldots, e_n\) are affinely independent, meaning \(e_1 - 0, \ldots, e_n - 0\) is linearly independent, which is the statement that a certain determinant is nonzero. Because that determinant depends continuously on the points and \(U\) is dense, we may choose \(v_0, v_1, \ldots, v_n \in U\) so close to \(0, e_1, \ldots, e_n\) respectively that \(v_1 - v_0, \ldots, v_n - v_0\) is still linearly independent. Convexity of \(U\) then implies that \(U\) contains the simplex

\begin{equation*} \Bigl\{ \sum_{j=0}^{n} t_j v_j : t_j \geq 0 \text{ for each } j \text{ and } \sum_{j=0}^{n} t_j = 1 \Bigr\}, \end{equation*}

which has nonempty interior (it is the image of the standard simplex under an invertible affine map). Thus \(U\) has nonempty interior. But then the nonempty open set \(\operatorname{int} U\) is disjoint from \(V \setminus U\), so \(V \setminus U\) is not dense in \(V\). Hence in finite dimensions the two conditions \(\overline{U} = V\) and \(\overline{V \setminus U} = V\) cannot both hold for a convex \(U\), even without requiring \(V \setminus U\) to be convex.

The essential point in the infinite-dimensional construction is the existence of a discontinuous linear functional (6.62), which has no finite-dimensional analogue: by Exercise 7 in Section 6D, every linear functional on a finite-dimensional normed vector space is continuous.

Exercises 6E

Problem (6E.1)

Suppose \(U\) is a subset of a metric space \(V\). Show that \(U\) is dense in \(V\) if and only if every nonempty open subset of \(V\) contains at least one element of \(U\).

Solution

Recall that \(U\) is dense in \(V\) means \(\overline{U} = V\) (6.75), and that by the definition of closure (6.7),

\begin{equation*} \overline{U} = \{ g \in V : B(g,\varepsilon) \cap U \neq \emptyset \text{ for every } \varepsilon > 0 \}. \end{equation*}

First suppose \(U\) is dense in \(V\), and let \(G\) be a nonempty open subset of \(V\). Choose \(f \in G\). Because \(G\) is open, the definition of open set (6.4) gives an \(r > 0\) such that \(B(f,r) \subseteq G\). Since \(\overline{U} = V\), we have \(f \in \overline{U}\), and hence \(B(f,r) \cap U \neq \emptyset\). Any element of \(B(f,r) \cap U\) is an element of \(U\) lying in \(G\). Thus every nonempty open subset of \(V\) contains at least one element of \(U\).

Conversely, suppose every nonempty open subset of \(V\) contains at least one element of \(U\). Let \(f \in V\) and let \(\varepsilon > 0\). The open ball \(B(f,\varepsilon)\) is an open subset of \(V\) (6.5), and it is nonempty because \(d(f,f) = 0 < \varepsilon\) shows \(f \in B(f,\varepsilon)\). Hence by hypothesis \(B(f,\varepsilon) \cap U \neq \emptyset\). As \(\varepsilon > 0\) was arbitrary, 6.7 gives \(f \in \overline{U}\). Because \(f \in V\) was arbitrary, \(V \subseteq \overline{U}\); the reverse inclusion holds trivially, so \(\overline{U} = V\). Thus \(U\) is dense in \(V\).

Problem (6E.2)

Suppose \(U\) is a subset of a metric space \(V\). Show that \(U\) has an empty interior if and only if \(V \setminus U\) is dense in \(V\).

Solution

Recall the definition of interior (6.74): \(f \in \operatorname{int} U\) means \(f \in U\) and \(B(f,r) \subseteq U\) for some \(r > 0\). We will use the characterization of density proved in Exercise 1: \(V \setminus U\) is dense in \(V\) if and only if every nonempty open subset of \(V\) contains at least one element of \(V \setminus U\).

First suppose \(\operatorname{int} U = \emptyset\). Let \(G\) be a nonempty open subset of \(V\); we must produce an element of \(G \cap (V \setminus U)\). Suppose no such element exists. Then \(G \subseteq U\). Choose \(f \in G\). Because \(G\) is open, 6.4 provides \(r > 0\) with \(B(f,r) \subseteq G \subseteq U\). Since also \(f \in G \subseteq U\), the definition 6.74 gives \(f \in \operatorname{int} U\), contradicting \(\operatorname{int} U = \emptyset\). Hence \(G \cap (V \setminus U) \neq \emptyset\). As \(G\) was an arbitrary nonempty open subset of \(V\), Exercise 1 implies that \(V \setminus U\) is dense in \(V\).

Conversely, suppose \(V \setminus U\) is dense in \(V\), and suppose (seeking a contradiction) that some \(f \in \operatorname{int} U\). Then there is \(r > 0\) with \(B(f,r) \subseteq U\). The set \(B(f,r)\) is open (6.5) and nonempty (it contains \(f\)), so by Exercise 1 it contains some element \(g \in V \setminus U\). But \(B(f,r) \subseteq U\) forces \(g \in U\), a contradiction. Hence \(\operatorname{int} U = \emptyset\).

Problem (6E.3)

Prove or give a counterexample: If \(V\) is a metric space and \(U, W\) are subsets of \(V\), then \((\operatorname{int} U) \cup (\operatorname{int} W) = \operatorname{int}(U \cup W)\).

Solution

The statement is false; only the inclusion \(\subseteq\) holds in general.

The inclusion that is always valid: if \(f \in \operatorname{int} U\), then \(B(f,r) \subseteq U \subseteq U \cup W\) for some \(r > 0\), so \(f \in \operatorname{int}(U \cup W)\) by 6.74; the same argument applies with \(W\) in place of \(U\). Hence

\begin{equation*} (\operatorname{int} U) \cup (\operatorname{int} W) \subseteq \operatorname{int}(U \cup W). \end{equation*}

Counterexample to the reverse inclusion. Let \(V = \mathbb{R}\) with its standard metric \(d(x,y) = |x-y|\), and let

\begin{equation*} U = [0,1], \qquad W = [1,2]. \end{equation*}

Then \(\operatorname{int} U = (0,1)\) and \(\operatorname{int} W = (1,2)\). Indeed, every open interval about a point of \((0,1)\) of small enough radius lies in \([0,1]\), while no open ball \(B(0,r) = (-r,r)\) is contained in \([0,1]\) (it contains \(-r/2 < 0\)) and no open ball \(B(1,r)\) is contained in \([0,1]\) (it contains \(1 + r/2 > 1\)); similarly for \(W\). Therefore

\begin{equation*} (\operatorname{int} U) \cup (\operatorname{int} W) = (0,1) \cup (1,2) = (0,2) \setminus \{1\}. \end{equation*}

On the other hand \(U \cup W = [0,2]\), and by the same reasoning as above \(\operatorname{int}[0,2] = (0,2)\). Thus

\begin{equation*} \operatorname{int}(U \cup W) = (0,2) \neq (0,2) \setminus \{1\} = (\operatorname{int} U) \cup (\operatorname{int} W), \end{equation*}

the point \(1\) lying in \(\operatorname{int}(U \cup W)\) but not in \((\operatorname{int} U) \cup (\operatorname{int} W)\).

An even more dramatic counterexample: take \(U = \mathbb{Q}\) and \(W = \mathbb{R} \setminus \mathbb{Q}\) in \(V = \mathbb{R}\). Both \(\mathbb{Q}\) and \(\mathbb{R} \setminus \mathbb{Q}\) are dense in \(\mathbb{R}\), so by Exercise 2 (applied to each of them, using that the complement of each is the other) both have empty interior. Hence \((\operatorname{int} U) \cup (\operatorname{int} W) = \emptyset\), while \(\operatorname{int}(U \cup W) = \operatorname{int} \mathbb{R} = \mathbb{R}\).

Problem (6E.4)

Prove or give a counterexample: If \(V\) is a metric space and \(U, W\) are subsets of \(V\), then \((\operatorname{int} U) \cap (\operatorname{int} W) = \operatorname{int}(U \cap W)\).

Solution

The statement is true. We prove both inclusions, using the definition of interior (6.74): \(f \in \operatorname{int} E\) means \(f \in E\) and \(B(f,r) \subseteq E\) for some \(r > 0\).

First suppose \(f \in (\operatorname{int} U) \cap (\operatorname{int} W)\). Then \(f \in U\) and \(f \in W\), so \(f \in U \cap W\). Also there exist \(r_1 > 0\) and \(r_2 > 0\) with \(B(f,r_1) \subseteq U\) and \(B(f,r_2) \subseteq W\). Let \(r = \min\{r_1, r_2\} > 0\). Because \(r \leq r_1\) and \(r \leq r_2\), we have \(B(f,r) \subseteq B(f,r_1) \subseteq U\) and \(B(f,r) \subseteq B(f,r_2) \subseteq W\), hence \(B(f,r) \subseteq U \cap W\). Thus \(f \in \operatorname{int}(U \cap W)\), proving

\begin{equation*} (\operatorname{int} U) \cap (\operatorname{int} W) \subseteq \operatorname{int}(U \cap W). \end{equation*}

Conversely, suppose \(f \in \operatorname{int}(U \cap W)\). Then \(f \in U \cap W\) and there exists \(r > 0\) with \(B(f,r) \subseteq U \cap W\). Since \(U \cap W \subseteq U\), we get \(f \in U\) and \(B(f,r) \subseteq U\), so \(f \in \operatorname{int} U\); since \(U \cap W \subseteq W\), we likewise get \(f \in \operatorname{int} W\). Hence \(f \in (\operatorname{int} U) \cap (\operatorname{int} W)\), proving the reverse inclusion.

Combining the two inclusions gives \((\operatorname{int} U) \cap (\operatorname{int} W) = \operatorname{int}(U \cap W)\).

[Conceptually: \(\operatorname{int} E\) is the largest open subset of \(V\) contained in \(E\), and the intersection of two open sets is open, which is exactly why the argument works here. For unions the corresponding statement fails, as Exercise 3 shows, because the union of the largest open set inside \(U\) and the largest open set inside \(W\) need not be the largest open set inside \(U \cup W\).]

Problem (6E.5)

Suppose

\begin{equation*} X = \{0\} \cup \bigcup_{k=1}^{\infty} \left\{ \tfrac{1}{k} \right\} \end{equation*}

and \(d(x,y) = |x-y|\) for \(x, y \in X\).

(a) Show that \((X,d)\) is a complete metric space.

(b) Each set of the form \(\{x\}\) for \(x \in X\) is a closed subset of \(\mathbb{R}\) that has an empty interior as a subset of \(\mathbb{R}\). Clearly \(X\) is a countable union of such sets. Explain why this does not violate the statement of Baire’s Theorem that a complete metric space is not the countable union of closed subsets with empty interior.

Solution

(a) Since \(d\) is the restriction to \(X\) of the standard metric on \(\mathbb{R}\), \((X,d)\) is a metric space. We show that \(X\) is a closed subset of \(\mathbb{R}\) and then deduce completeness from the completeness of \(\mathbb{R}\).

\(X\) is closed in \(\mathbb{R}\): we check that \(\mathbb{R} \setminus X\) is open by exhibiting it as a union of open intervals. Note \(X \subseteq [0,1]\). If \(x \in \mathbb{R} \setminus X\), then exactly one of the following holds: \(x < 0\); \(x > 1\); or \(0 < x < 1\) and \(x\) is not of the form \(1/k\). In the last case, since \(1/k \to 0\) as \(k \to \infty\) and \(1/1 = 1\), there is a unique positive integer \(k\) with \(\frac{1}{k+1} < x < \frac{1}{k}\). Hence

\begin{equation*} \mathbb{R} \setminus X = (-\infty, 0) \cup (1, \infty) \cup \bigcup_{k=1}^{\infty} \left( \tfrac{1}{k+1}, \tfrac{1}{k} \right). \end{equation*}

Each set on the right is an open subset of \(\mathbb{R}\), and a union of open sets is open. Thus \(\mathbb{R} \setminus X\) is open, so \(X\) is closed in \(\mathbb{R}\) by 6.6.

Now suppose \(x_1, x_2, \ldots\) is a Cauchy sequence in \(X\). It is then a Cauchy sequence in \(\mathbb{R}\), and \(\mathbb{R}\) with its standard metric is complete, so there exists \(x \in \mathbb{R}\) with \(\lim_{k \to \infty} x_k = x\). Because \(X\) is a closed subset of \(\mathbb{R}\) and \(x\) is the limit of a convergent sequence of elements of \(X\), we have \(x \in X\) by 6.9(e). Since \(d\) is the restriction of the metric of \(\mathbb{R}\), the statement \(\lim_{k \to \infty} d(x_k, x) = 0\) says precisely that \(x_k \to x\) in \((X,d)\). Thus every Cauchy sequence in \(X\) converges to an element of \(X\), which means \((X,d)\) is a complete metric space.

(b) The hypothesis in Baire’s Theorem [6.76(a)] is about closed subsets that have empty interior as subsets of the complete metric space in question, and here that space is \(X\), not \(\mathbb{R}\). The sets \(\{x\}\) for \(x \in X\) do have empty interior as subsets of \(\mathbb{R}\), but almost all of them fail to have empty interior as subsets of \(X\), so the decomposition \(X = \bigcup_{x \in X} \{x\}\) is not a decomposition of the kind that Baire’s Theorem forbids.

Concretely, fix \(k \in \mathbb{Z}^+\) and let \(r = \frac{1}{k(k+1)} > 0\). The elements of \(X\) other than \(\frac{1}{k}\) are \(0\) and the numbers \(\frac{1}{j}\) with \(j \neq k\). We have

\begin{equation*} \left| \tfrac{1}{k} - 0 \right| = \tfrac{1}{k} > \tfrac{1}{k(k+1)} = r, \end{equation*}

and if \(j > k\) then \(\left| \frac{1}{k} - \frac{1}{j} \right| \geq \frac{1}{k} - \frac{1}{k+1} = \frac{1}{k(k+1)} = r\), while if \(j < k\) then \(\left| \frac{1}{k} - \frac{1}{j} \right| \geq \frac{1}{k-1} - \frac{1}{k} = \frac{1}{k(k-1)} > r\). Hence the open ball of \(X\) centered at \(\frac{1}{k}\) with radius \(r\) is

\begin{equation*} B\!\left( \tfrac{1}{k}, r \right) = \left\{ y \in X : \left| y - \tfrac{1}{k} \right| < r \right\} = \left\{ \tfrac{1}{k} \right\}. \end{equation*}

So \(B\left(\frac{1}{k}, r\right) \subseteq \left\{\frac{1}{k}\right\}\), which by 6.74 means \(\operatorname{int}\left\{\frac{1}{k}\right\} = \left\{\frac{1}{k}\right\} \neq \emptyset\) in \(X\). Every singleton \(\left\{\frac{1}{k}\right\}\) is therefore an open subset of \(X\) with nonempty interior in \(X\).

The only element of \(X\) whose singleton does have empty interior in \(X\) is \(0\): for every \(r > 0\) there is \(k \in \mathbb{Z}^+\) with \(\frac{1}{k} < r\), so \(B(0,r)\) contains the point \(\frac{1}{k} \neq 0\) and hence \(B(0,r) \not\subseteq \{0\}\).

Thus in the countable union \(X = \{0\} \cup \bigcup_{k=1}^{\infty} \left\{\frac{1}{k}\right\}\), all but one of the closed sets involved have nonempty interior in \(X\), and there is no conflict with 6.76(a). (Indeed 6.76(a) guarantees that \(X\), being complete, cannot be written as a countable union of closed subsets of \(X\) each having empty interior in \(X\), and the displayed decomposition is not such a union.)

Problem (6E.6)

Give an example of a metric space that is the countable union of closed subsets with empty interior.

[This exercise shows that the completeness hypothesis in Baire’s Theorem cannot be dropped.]

Solution

Take \(V = \mathbb{Q}\) with the metric \(d(x,y) = |x-y|\) inherited from \(\mathbb{R}\). This is a metric space, and it is countable, so we may enumerate its elements as \(\mathbb{Q} = \{r_1, r_2, \ldots\}\). Then

\begin{equation*} \mathbb{Q} = \bigcup_{j=1}^{\infty} \{ r_j \}, \end{equation*}

a countable union. We claim each \(\{r_j\}\) is a closed subset of \(\mathbb{Q}\) with empty interior in \(\mathbb{Q}\).

\(\{r_j\}\) is closed in \(\mathbb{Q}\): its complement \(\mathbb{Q} \setminus \{r_j\}\) is open, because if \(x \in \mathbb{Q}\) and \(x \neq r_j\), then \(r = |x - r_j| > 0\) and \(B(x,r) \subseteq \mathbb{Q} \setminus \{r_j\}\) (every \(y\) with \(|y - x| < r\) satisfies \(y \neq r_j\)). Hence \(\{r_j\}\) is closed by 6.6.

\(\{r_j\}\) has empty interior in \(\mathbb{Q}\): let \(r > 0\). Since \(\mathbb{Q}\) is dense in \(\mathbb{R}\), the interval \((r_j, r_j + r)\) contains a rational number \(q\); then \(q \in \mathbb{Q}\), \(|q - r_j| < r\), and \(q \neq r_j\). Thus \(B(r_j, r) \not\subseteq \{r_j\}\) for every \(r > 0\), so by 6.74 no point of \(\{r_j\}\) is an interior point, i.e. \(\operatorname{int}\{r_j\} = \emptyset\).

Hence \(\mathbb{Q}\) is a countable union of closed subsets of \(\mathbb{Q}\) each having empty interior in \(\mathbb{Q}\).

By 6.76(a) this forces \(\mathbb{Q}\) to be incomplete, and indeed it is: choose a sequence of rationals converging in \(\mathbb{R}\) to \(\sqrt{2}\) (possible because \(\mathbb{Q}\) is dense in \(\mathbb{R}\): for each \(k\) pick \(x_k \in \mathbb{Q}\) with \(|x_k - \sqrt{2}| < \frac{1}{k}\)). That sequence is Cauchy in \(\mathbb{Q}\), since it converges in \(\mathbb{R}\), but it has no limit in \(\mathbb{Q}\): a limit in \(\mathbb{Q}\) would also be a limit in \(\mathbb{R}\), and limits in \(\mathbb{R}\) are unique, so the limit would equal the irrational number \(\sqrt{2}\), which is not in \(\mathbb{Q}\).

So the completeness hypothesis in 6.76(a) cannot be dropped.

[The same argument works for any countable metric space in which no singleton is open, for instance \(\mathbb{Q} \cap (0,1)\) or the set of dyadic rationals, each with the metric inherited from \(\mathbb{R}\). Contrast this with Exercise 5, where \(X\) is countable but almost every singleton is open in \(X\), so the corresponding decomposition of \(X\) is harmless.]

Problem (6E.7)

(a) Define \(f : \mathbb{R} \to \mathbb{R}\) as follows:

\begin{equation*} f(a) = \begin{cases} 0 & \text{if } a \text{ is irrational,} \\ \frac{1}{n} & \text{if } a \text{ is rational and } n \text{ is the smallest positive integer such that } a = \frac{m}{n} \text{ for some integer } m. \end{cases} \end{equation*}

At which numbers in \(\mathbb{R}\) is \(f\) continuous?

(b) Show that there does not exist a countable collection of open subsets of \(\mathbb{R}\) whose intersection equals \(\mathbb{Q}\).

(c) Show that there does not exist a function \(f : \mathbb{R} \to \mathbb{R}\) such that \(f\) is continuous at each element of \(\mathbb{Q}\) and discontinuous at each element of \(\mathbb{R} \setminus \mathbb{Q}\).

Solution

(a) The function \(f\) is continuous at every irrational number and discontinuous at every rational number.

First note that \(f\) is well defined: for \(a \in \mathbb{Q}\) the set of positive integers \(n\) such that \(a = m/n\) for some \(m \in \mathbb{Z}\) is nonempty, hence has a smallest element. Write \(n(a)\) for that smallest \(n\), so \(f(a) = 1/n(a) > 0\) for every rational \(a\), and \(f(a) = 0\) for every irrational \(a\). Also \(0 \leq f \leq 1\) everywhere.

Throughout we use the definition of continuity at a point (6.10): \(f\) is continuous at \(a\) means that for every \(\varepsilon > 0\) there exists \(\delta > 0\) such that \(|f(x) - f(a)| < \varepsilon\) for all \(x \in \mathbb{R}\) with \(|x - a| < \delta\).

Discontinuity at each rational. Let \(a \in \mathbb{Q}\) and take \(\varepsilon = f(a) = \frac{1}{n(a)}\), which is positive. Let \(\delta > 0\) be arbitrary. Because \(\mathbb{R} \setminus \mathbb{Q}\) is dense in \(\mathbb{R}\), the interval \((a - \delta, a + \delta)\) contains some irrational \(x\); for that \(x\) we have \(|x - a| < \delta\) yet

\begin{equation*} |f(x) - f(a)| = |0 - f(a)| = f(a) = \varepsilon, \end{equation*}

so the required inequality \(|f(x) - f(a)| < \varepsilon\) fails. Since no \(\delta > 0\) works for this \(\varepsilon\), the function \(f\) is not continuous at \(a\).

Continuity at each irrational. Let \(a \in \mathbb{R} \setminus \mathbb{Q}\), so \(f(a) = 0\), and let \(\varepsilon > 0\). Choose \(N \in \mathbb{Z}^+\) with \(\frac{1}{N} < \varepsilon\). Let

\begin{equation*} S = \left\{ x \in (a-1, a+1) : f(x) \geq \varepsilon \right\}. \end{equation*}

If \(x \in S\) then \(f(x) > 0\), so \(x\) is rational, and \(\frac{1}{n(x)} = f(x) \geq \varepsilon > \frac{1}{N}\), hence \(n(x) < N\). Writing \(x = m/n(x)\), we conclude

\begin{equation*} S \subseteq \left\{ \tfrac{m}{n} : n \in \{1, \ldots, N\}, \ m \in \mathbb{Z}, \ \tfrac{m}{n} \in (a-1, a+1) \right\}. \end{equation*}

For each fixed \(n \leq N\) the condition \(\frac{m}{n} \in (a-1, a+1)\) means \(m \in (n(a-1), n(a+1))\), an interval of length \(2n \leq 2N\), which contains at most \(2N + 1\) integers. Hence \(S\) is a finite set. Moreover \(a \notin S\), because \(a\) is irrational and every element of \(S\) is rational.

If \(S = \emptyset\), put \(\delta = 1\); otherwise put

\begin{equation*} \delta = \min \left( \{1\} \cup \{ |a - s| : s \in S \} \right), \end{equation*}

which is positive because \(S\) is finite and \(a \notin S\). Now suppose \(|x - a| < \delta\). Then \(x \in (a-1, a+1)\) and \(x \notin S\), so by the definition of \(S\) we have \(f(x) < \varepsilon\). Since \(f \geq 0\) and \(f(a) = 0\),

\begin{equation*} |f(x) - f(a)| = f(x) < \varepsilon. \end{equation*}

Thus \(f\) is continuous at \(a\).

Therefore the set of points of continuity of \(f\) is exactly \(\mathbb{R} \setminus \mathbb{Q}\).

(b) Suppose, seeking a contradiction, that \(G_1, G_2, \ldots\) are open subsets of \(\mathbb{R}\) with

\begin{equation*} \bigcap_{k=1}^{\infty} G_k = \mathbb{Q}. \end{equation*}

Then \(\mathbb{Q} \subseteq G_k\) for every \(k\), so \(\mathbb{R} = \overline{\mathbb{Q}} \subseteq \overline{G_k} \subseteq \mathbb{R}\), which shows that each \(G_k\) is a dense open subset of \(\mathbb{R}\).

Enumerate the rationals as \(\mathbb{Q} = \{r_1, r_2, \ldots\}\), and for each \(j \in \mathbb{Z}^+\) set

\begin{equation*} H_j = \mathbb{R} \setminus \{ r_j \}. \end{equation*}

Each \(H_j\) is open (if \(x \neq r_j\) then \(B(x, |x - r_j|) \subseteq H_j\)) and each \(H_j\) is dense in \(\mathbb{R}\): given any nonempty open \(G \subseteq \mathbb{R}\), choose \(x \in G\) and \(r > 0\) with \(B(x,r) \subseteq G\); the interval \(B(x,r)\) is infinite, so it contains a point other than \(r_j\), and that point lies in \(G \cap H_j\). By Exercise 1 this gives density.

Now \(G_1, H_1, G_2, H_2, \ldots\) is a countable collection of dense open subsets of \(\mathbb{R}\), and

\begin{equation*} \left( \bigcap_{k=1}^{\infty} G_k \right) \cap \left( \bigcap_{j=1}^{\infty} H_j \right) = \mathbb{Q} \cap \left( \mathbb{R} \setminus \mathbb{Q} \right) = \emptyset, \end{equation*}

because \(\bigcap_{j=1}^{\infty} H_j = \mathbb{R} \setminus \{r_1, r_2, \ldots\} = \mathbb{R} \setminus \mathbb{Q}\). But \(\mathbb{R}\) with its standard metric is a complete metric space, so Baire’s Theorem [6.76(b)] says that the countable intersection of dense open subsets of \(\mathbb{R}\) is nonempty. This contradiction shows that no such collection \(G_1, G_2, \ldots\) exists.

(c) Suppose \(f : \mathbb{R} \to \mathbb{R}\) is any function. For \(a \in \mathbb{R}\) and \(\delta > 0\) define the oscillation

\begin{equation*} \omega(a, \delta) = \sup \left\{ |f(x) - f(y)| : x, y \in (a - \delta, a + \delta) \right\} \in [0, \infty], \end{equation*}

and set \(\omega(a) = \inf_{\delta > 0} \omega(a, \delta)\). Note that \(\omega(a, \delta)\) is a nondecreasing function of \(\delta\).

Claim: \(f\) is continuous at \(a\) if and only if \(\omega(a) = 0\).

If \(f\) is continuous at \(a\) and \(\varepsilon > 0\), then 6.10 provides \(\delta > 0\) such that \(|f(x) - f(a)| < \varepsilon/3\) whenever \(|x - a| < \delta\). Then for \(x, y \in (a - \delta, a + \delta)\) we have \(|f(x) - f(y)| \leq |f(x) - f(a)| + |f(a) - f(y)| < 2\varepsilon/3\), so \(\omega(a, \delta) \leq 2\varepsilon/3 < \varepsilon\) and hence \(\omega(a) < \varepsilon\). As \(\varepsilon > 0\) was arbitrary, \(\omega(a) = 0\). Conversely, if \(\omega(a) = 0\) and \(\varepsilon > 0\), there is \(\delta > 0\) with \(\omega(a, \delta) < \varepsilon\); then \(|x - a| < \delta\) gives \(|f(x) - f(a)| \leq \omega(a, \delta) < \varepsilon\) (take \(y = a\) in the supremum). So \(f\) is continuous at \(a\). This proves the claim.

Now for \(n \in \mathbb{Z}^+\) let

\begin{equation*} G_n = \left\{ a \in \mathbb{R} : \omega(a) < \tfrac{1}{n} \right\}. \end{equation*}

Each \(G_n\) is open: if \(\omega(a) < \frac{1}{n}\), choose \(\delta > 0\) with \(\omega(a, \delta) < \frac{1}{n}\). If \(|b - a| < \delta/2\), then \((b - \delta/2, b + \delta/2) \subseteq (a - \delta, a + \delta)\), so the supremum defining \(\omega(b, \delta/2)\) is over a subset of the pairs used for \(\omega(a, \delta)\), giving \(\omega(b) \leq \omega(b, \delta/2) \leq \omega(a, \delta) < \frac{1}{n}\) and hence \(b \in G_n\). Thus \(B(a, \delta/2) \subseteq G_n\), so \(G_n\) is open by 6.4.

By the claim, the set \(C\) of points at which \(f\) is continuous satisfies

\begin{equation*} C = \left\{ a \in \mathbb{R} : \omega(a) = 0 \right\} = \bigcap_{n=1}^{\infty} G_n, \end{equation*}

a countable intersection of open subsets of \(\mathbb{R}\).

Finally, suppose there were a function \(f : \mathbb{R} \to \mathbb{R}\) continuous at each element of \(\mathbb{Q}\) and discontinuous at each element of \(\mathbb{R} \setminus \mathbb{Q}\). Then its set of continuity points is \(C = \mathbb{Q}\), so by the previous paragraph \(\mathbb{Q}\) would be the intersection of the countable collection \(G_1, G_2, \ldots\) of open subsets of \(\mathbb{R}\). This contradicts part (b). Hence no such function exists.

[Compare part (a): there the roles are reversed, and indeed the set of continuity points \(\mathbb{R} \setminus \mathbb{Q} = \bigcap_{j=1}^{\infty} \left( \mathbb{R} \setminus \{r_j\} \right)\) is a countable intersection of open sets, consistent with the description of \(C\) above.]

Problem (6E.8)

Suppose \((X, d)\) is a complete metric space and \(G_1, G_2, \ldots\) is a sequence of dense open subsets of \(X\). Prove that \(\bigcap_{k=1}^{\infty} G_k\) is a dense subset of \(X\).

Solution

Throughout, \(B(f, r) = \{x \in X : d(x, f) < r\}\) denotes the open ball and \(\bar{B}(f, r) = \{x \in X : d(x, f) \le r\}\) denotes the closed ball.

We use three preliminary facts.

Fact 1: if \(G\) is an open subset of a metric space and \(f \in G\), then there exists \(r > 0\) with \(\overline{B(f, r)} \subseteq G\). This is the fact whose proof is requested just before 6.76. Indeed, by the definition of an open set there exists \(s > 0\) with \(B(f, s) \subseteq G\); take \(r = s/2\). The closed ball \(\bar{B}(f, r)\) is a closed set (its complement \(\{x : d(x, f) > r\}\) is open, because if \(d(x, f) > r\) then \(B(x, d(x,f) - r)\) misses \(\bar{B}(f,r)\) by the triangle inequality) containing \(B(f, r)\), so \(\overline{B(f, r)} \subseteq \bar{B}(f, r) \subseteq B(f, s) \subseteq G\).

Fact 2: a closed subset \(Y\) of a complete metric space \((X, d)\) is itself a complete metric space under the restricted metric. This is 6.16(b). (Directly: a Cauchy sequence in \(Y\) is a Cauchy sequence in \(X\), hence converges to some \(g \in X\); because \(Y\) is closed and the sequence lies in \(Y\), we have \(g \in Y\) by 6.9(e), and the convergence takes place in \(Y\).)

Fact 3: a subset \(U\) of a metric space \(X\) is dense in \(X\) if and only if every nonempty open subset of \(X\) contains an element of \(U\). This is the first elementary fact listed after 6.75 (Exercise 1 in Section 6E).

Now let \(U\) be an arbitrary nonempty open subset of \(X\). By Fact 3, it suffices to show that \(U \cap \bigcap_{k=1}^{\infty} G_k \ne \emptyset\).

Choose \(f \in U\). By Fact 1 there exists \(r > 0\) with \(\overline{B(f, r)} \subseteq U\). Let

\begin{equation*} Y = \overline{B(f, r)}, \end{equation*}

the closure in \(X\) of the open ball \(B(f, r)\). Thus

\begin{equation*} B(f, r) \subseteq Y \subseteq U . \end{equation*}

In particular \(Y\) is nonempty (it contains \(f\), since \(d(f,f) = 0 < r\)) and \(Y\) is a closed subset of \(X\) by 6.9(c); hence by Fact 2 the metric space \((Y, d)\) is complete.

For each \(k \in \mathbf{Z}^+\) set \(H_k = G_k \cap Y\). We claim each \(H_k\) is a dense open subset of the metric space \(Y\).

\(H_k\) is open in \(Y\): if \(y \in H_k\), then \(y \in G_k\) and \(G_k\) is open in \(X\), so there is \(\rho > 0\) with \(B(y, \rho) \subseteq G_k\); hence the ball of \(Y\) centered at \(y\) of radius \(\rho\), which equals \(B(y, \rho) \cap Y\), is contained in \(G_k \cap Y = H_k\).

\(H_k\) is dense in \(Y\): let \(y \in Y\) and \(\varepsilon > 0\). Because \(y\) lies in the closure of \(B(f, r)\), there exists \(x \in B(f, r)\) with \(d(x, y) < \varepsilon/2\). The set \(B(f, r) \cap B(x, \varepsilon/2)\) is an open subset of \(X\) and it is nonempty (it contains \(x\)). Because \(G_k\) is dense in \(X\), Fact 3 provides

\begin{equation*} z \in G_k \cap B(f, r) \cap B(x, \tfrac{\varepsilon}{2}). \end{equation*}

Then \(z \in G_k \cap B(f, r) \subseteq G_k \cap Y = H_k\) and

\begin{equation*} d(z, y) \le d(z, x) + d(x, y) < \tfrac{\varepsilon}{2} + \tfrac{\varepsilon}{2} = \varepsilon . \end{equation*}

Thus every open ball of \(Y\) centered at a point of \(Y\) contains an element of \(H_k\). Because every nonempty open subset of \(Y\) contains such a ball, Fact 3 (applied in the metric space \(Y\)) shows that \(H_k\) is dense in \(Y\).

We have shown that \(H_1, H_2, \ldots\) is a sequence of dense open subsets of the complete metric space \((Y, d)\). By Baire’s Theorem [6.76(b)] there exists

\begin{equation*} g \in \bigcap_{k=1}^{\infty} H_k = Y \cap \bigcap_{k=1}^{\infty} G_k . \end{equation*}

Because \(Y \subseteq U\), this \(g\) belongs to \(U \cap \bigcap_{k=1}^{\infty} G_k\), so that intersection is nonempty.

Since \(U\) was an arbitrary nonempty open subset of \(X\), Fact 3 shows that \(\bigcap_{k=1}^{\infty} G_k\) is dense in \(X\).

Problem (6E.9)

Prove that there does not exist an infinite-dimensional Banach space with a countable basis.

[This exercise implies, for example, that there is not a norm that makes the vector space of polynomials with coefficients in \(\mathbf{F}\) into a Banach space.]

Solution

Here basis has the meaning given in 6.54: a linearly independent family whose span is the whole space, where the span consists of the finite linear combinations.

Suppose, seeking a contradiction, that \(V\) is an infinite-dimensional Banach space with a countable basis \(\{e_k\}_{k \in \Gamma}\). If \(\Gamma\) were finite then \(V = \operatorname{span}\{e_k\}_{k \in \Gamma}\) would be finite-dimensional by 6.54, contradicting our hypothesis. Hence \(\Gamma\) is countably infinite, and after relabeling we may assume \(\Gamma = \mathbf{Z}^+\). So \(e_1, e_2, \ldots\) is a basis of \(V\).

For \(n \in \mathbf{Z}^+\) let

\begin{equation*} V_n = \operatorname{span}\{e_1, \ldots, e_n\}. \end{equation*}

Step 1: \(V = \bigcup_{n=1}^{\infty} V_n\). Each \(V_n\) is a subset of \(V\), giving one inclusion. Conversely, if \(f \in V\) then, because the span of \(e_1, e_2, \ldots\) equals \(V\), there exist a finite set \(\Omega \subseteq \mathbf{Z}^+\) and scalars \(\{\alpha_j\}_{j \in \Omega}\) with \(f = \sum_{j \in \Omega} \alpha_j e_j\). Choosing \(n = \max \Omega\) (or \(n = 1\) if \(\Omega = \emptyset\)) gives \(f \in V_n\).

Step 2: each \(V_n\) is a closed subset of \(V\). This is exactly Exercise 9 in Section 6D: every finite-dimensional subspace of a normed vector space is closed. (Briefly: \(V_n\), with the norm inherited from \(V\), is a finite-dimensional normed vector space, hence a Banach space by Exercise 8 in Section 6D; thus \(V_n\) is a complete subset of the metric space \(V\), and 6.16(a) shows that a complete subset of a metric space is closed.)

Step 3: each \(V_n\) has empty interior in \(V\). Suppose instead that some \(V_n\) has nonempty interior. By the definition of interior (6.74) there then exist \(g \in V_n\) and \(\rho > 0\) with \(B(g, \rho) \subseteq V_n\). (This is the situation of Exercise 10 in Section 6C, whose conclusion is that \(V_n = V\); here is the argument.) Let \(f \in V\) with \(f \ne 0\). Then

\begin{equation*} \left\| \left( g + \frac{\rho}{2 \|f\|} f \right) - g \right\| = \frac{\rho}{2} < \rho , \end{equation*}

so \(g + \frac{\rho}{2 \|f\|} f \in B(g, \rho) \subseteq V_n\). Because \(V_n\) is a subspace containing \(g\), we conclude that \(\frac{\rho}{2\|f\|} f \in V_n\) and hence \(f \in V_n\). As \(0 \in V_n\) trivially, this shows \(V = V_n\), which makes \(V\) finite-dimensional, contradicting our hypothesis. Hence \(\operatorname{int} V_n = \emptyset\).

Now \(V\), with the metric \(d(f, g) = \|f - g\|\), is a complete metric space because \(V\) is a Banach space (6.37). Steps 1, 2, 3 exhibit \(V\) as the countable union of closed subsets with empty interior, which contradicts Baire’s Theorem [6.76(a)].

This contradiction shows that no infinite-dimensional Banach space has a countable basis.

For the bracketed remark: let \(P\) denote the vector space of all polynomials with coefficients in \(\mathbf{F}\). The family \(1, x, x^2, \ldots\) is linearly independent and spans \(P\), so it is a countable basis of \(P\), and \(P\) is infinite-dimensional (no finite family spans \(P\), since the span of finitely many polynomials contains only polynomials of degree at most the largest degree occurring). Hence by what we proved, no norm on \(P\) makes \(P\) into a Banach space.

Problem (6E.10)

Give an example of a Banach space \(V\), a normed vector space \(W\), a bounded linear map \(T\) of \(V\) onto \(W\), and an open subset \(G\) of \(V\) such that \(T(G)\) is not an open subset of \(W\).

[This exercise shows that the hypothesis in the Open Mapping Theorem that \(W\) is a Banach space cannot be relaxed to the hypothesis that \(W\) is a normed vector space.]

Solution

Let \(\ell^1\) denote the vector space of sequences \((a_1, a_2, \ldots)\) of elements of \(\mathbf{F}\) with \(\sum_{k=1}^{\infty} |a_k| < \infty\). Take

\begin{equation*} V = \ell^1 \text{ with the norm } \|(a_1, a_2, \ldots)\|_1 = \sum_{k=1}^{\infty} |a_k| , \end{equation*}

\begin{equation*} W = \ell^1 \text{ with the norm } \|(a_1, a_2, \ldots)\|_{\infty} = \sup_{k \in \mathbf{Z}^+} |a_k| , \end{equation*}

and let \(T : V \to W\) be the identity map \(Ta = a\). Finally let

\begin{equation*} G = \{a \in \ell^1 : \|a\|_1 < 1\}, \end{equation*}

the open unit ball of \(V\), which is an open subset of \(V\).

By the second bullet point of Example 6.38, \(V\) is a Banach space. By the second bullet point of Example 6.39, \(W\) is a normed vector space that is not a Banach space.

\(T\) is clearly linear, and it maps \(V\) onto \(W\) because it is the identity map on the common underlying vector space \(\ell^1\) (in particular \(T\) is also one-to-one). \(T\) is bounded because for every \(a = (a_1, a_2, \ldots) \in \ell^1\) and every \(k \in \mathbf{Z}^+\) we have \(|a_k| \le \sum_{j=1}^{\infty} |a_j|\), and hence

\begin{equation*} \|Ta\|_{\infty} = \sup_{k \in \mathbf{Z}^+} |a_k| \le \|a\|_1 . \end{equation*}

Thus \(\|T\| \le 1\).

Finally we show that \(T(G) = \{a \in \ell^1 : \|a\|_1 < 1\}\) is not an open subset of \(W\). Note \(0 \in T(G)\). Suppose there existed \(\varepsilon > 0\) with

\begin{equation*} \{a \in \ell^1 : \|a\|_{\infty} < \varepsilon\} \subseteq T(G). \end{equation*}

Choose \(n \in \mathbf{Z}^+\) with \(n > 2/\varepsilon\) and let

\begin{equation*} a = \Big( \underbrace{\tfrac{\varepsilon}{2}, \ldots, \tfrac{\varepsilon}{2}}_{n \text{ terms}}, 0, 0, \ldots \Big). \end{equation*}

Then \(a \in \ell^1\) and \(\|a\|_{\infty} = \varepsilon/2 < \varepsilon\), so \(a\) would lie in \(T(G)\). But

\begin{equation*} \|a\|_1 = n \cdot \frac{\varepsilon}{2} > \frac{2}{\varepsilon} \cdot \frac{\varepsilon}{2} = 1 , \end{equation*}

so \(a \notin T(G)\). This contradiction shows that no open ball of \(W\) centered at \(0\) is contained in \(T(G)\). Hence \(T(G)\) is not an open subset of \(W\).

There is no conflict with the Open Mapping Theorem (6.81) because \(W\) is not a Banach space.

Problem (6E.11)

Show that there exists a normed vector space \(V\), a Banach space \(W\), a bounded linear map \(T\) of \(V\) onto \(W\), and an open subset \(G\) of \(V\) such that \(T(G)\) is not an open subset of \(W\).

[This exercise shows that the hypothesis in the Open Mapping Theorem that \(V\) is a Banach space cannot be relaxed to the hypothesis that \(V\) is a normed vector space.]

Solution

Let \(W\) be any infinite-dimensional Banach space; to be concrete, let \(W = \ell^1\) with the norm \(\|(a_1, a_2, \ldots)\|_1 = \sum_{k=1}^{\infty} |a_k|\), which is a Banach space by the second bullet point of Example 6.38 and is infinite-dimensional. [For the last assertion, let \(e_k \in \ell^1\) denote the sequence with \(1\) in slot \(k\) and \(0\) in all other slots; clearly \(e_1, e_2, \ldots\) is a linearly independent family in \(\ell^1\). If \(\ell^1\) were finite-dimensional, then by 6.54 there would exist \(n \in \mathbf{Z}^+\) and \(f_1, \ldots, f_n \in \ell^1\) with \(\ell^1 = \operatorname{span}\{f_1, \ldots, f_n\}\). The standard linear algebra fact that in a vector space the length of every linearly independent list is at most the length of every spanning list would then force \(n + 1 \le n\), because \(e_1, \ldots, e_{n+1}\) is a linearly independent list of length \(n+1\). This contradiction shows \(\ell^1\) is infinite-dimensional.]

Because \(W\) is an infinite-dimensional normed vector space, 6.62 provides a discontinuous linear functional \(\varphi : W \to \mathbf{F}\). By 6.48, discontinuity means precisely that \(\varphi\) is not bounded; that is, we have the following property, to which we refer as the unboundedness of \(\varphi\):

\begin{equation*} \text{there is no } c \in (0, \infty) \text{ such that } |\varphi(f)| \le c \|f\|_1 \text{ for all } f \in W. \end{equation*}

Let \(V\) be the vector space \(\ell^1\) equipped with

\begin{equation*} \|f\|_V = \|f\|_1 + |\varphi(f)| . \end{equation*}

This is a norm on \(\ell^1\): homogeneity holds because \(\|\alpha f\|_1 + |\varphi(\alpha f)| = |\alpha| (\|f\|_1 + |\varphi(f)|)\) for \(\alpha \in \mathbf{F}\), the triangle inequality follows by adding the triangle inequalities for \(\|\cdot\|_1\) and for \(|\varphi(\cdot)|\) (the latter using linearity of \(\varphi\)), and \(\|f\|_V = 0\) forces \(\|f\|_1 = 0\), hence \(f = 0\). Thus \(V\) is a normed vector space.

Let \(T : V \to W\) be the identity map \(Tf = f\). Then \(T\) is linear and maps \(V\) onto \(W\) (indeed one-to-one and onto, since \(V\) and \(W\) have the same underlying vector space). Also \(T\) is bounded, because

\begin{equation*} \|Tf\|_1 = \|f\|_1 \le \|f\|_1 + |\varphi(f)| = \|f\|_V \end{equation*}

for all \(f \in V\); thus \(\|T\| \le 1\).

Let \(G = \{f \in V : \|f\|_V < 1\}\), the open unit ball of \(V\), which is an open subset of \(V\). Then

\begin{equation*} T(G) = \{f \in \ell^1 : \|f\|_1 + |\varphi(f)| < 1\}, \end{equation*}

and \(0 \in T(G)\).

Suppose \(T(G)\) were an open subset of \(W\). Then there would exist \(\varepsilon > 0\) with \(\{f \in W : \|f\|_1 < \varepsilon\} \subseteq T(G)\), and in particular

\begin{equation*} \|f\|_1 < \varepsilon \implies |\varphi(f)| < 1 . \end{equation*}

Now let \(f \in W\) with \(f \ne 0\), and set \(g = \frac{\varepsilon}{2\|f\|_1} f\). Then \(\|g\|_1 = \varepsilon/2 < \varepsilon\), so \(|\varphi(g)| < 1\), which by the homogeneity of \(\varphi\) means

\begin{equation*} |\varphi(f)| < \frac{2}{\varepsilon} \|f\|_1 . \end{equation*}

This inequality also holds (with \(\le\)) for \(f = 0\). Hence \(|\varphi(f)| \le \frac{2}{\varepsilon}\|f\|_1\) for all \(f \in W\), contradicting the unboundedness of \(\varphi\).

Therefore \(T(G)\) is not an open subset of \(W\), as desired.

No contradiction with the Open Mapping Theorem (6.81) arises: \(V\) is not a Banach space. (Indeed, were \(V\) complete, 6.81 applied to \(T : V \to W\) would make \(T(G)\) open.)

Problem (6E.12)

Suppose \(T : V \to W\) is a bounded linear map from a Banach space \(V\) to a Banach space \(W\). Prove that \(T\) is bounded below if and only if \(T\) is injective and the range of \(T\) is a closed subspace of \(W\).

Solution

Recall the definition stated just before this exercise: \(T\) is bounded below if there exists \(c \in (0, \infty)\) such that \(\|f\| \le c \|Tf\|\) for all \(f \in V\).

Write \(R = T(V)\) for the range of \(T\); because \(T\) is linear, \(R\) is a subspace of \(W\).

First suppose \(T\) is bounded below, say \(\|f\| \le c\|Tf\|\) for all \(f \in V\), where \(c \in (0, \infty)\).

\(T\) is injective: if \(Tf = 0\), then \(\|f\| \le c\|Tf\| = 0\), so \(f = 0\); because \(T\) is linear this implies \(T\) is one-to-one (if \(Tf = Tg\) then \(T(f - g) = 0\), so \(f = g\)).

\(R\) is closed: suppose \(g_1, g_2, \ldots\) is a sequence in \(R\) converging in \(W\) to some \(g \in W\). Write \(g_k = Tf_k\) with \(f_k \in V\). Because the convergent sequence \(g_1, g_2, \ldots\) is a Cauchy sequence (6.13) and

\begin{equation*} \|f_j - f_k\| \le c \|T(f_j - f_k)\| = c \|g_j - g_k\| \end{equation*}

for all \(j, k \in \mathbf{Z}^+\), the sequence \(f_1, f_2, \ldots\) is a Cauchy sequence in \(V\). Because \(V\) is a Banach space, there exists \(f \in V\) with \(\lim_{k \to \infty} f_k = f\). Because \(T\) is bounded, it is continuous (6.48), so

\begin{equation*} g = \lim_{k \to \infty} Tf_k = Tf \in R . \end{equation*}

Thus \(R\) contains the limit of every convergent sequence of elements of \(R\), which by 6.9(e) means \(R\) is a closed subset of \(W\). Hence the range of \(T\) is a closed subspace of \(W\).

Conversely, suppose \(T\) is injective and \(R\) is a closed subspace of \(W\).

Because \(W\) is a Banach space and \(R\) is a closed subset of \(W\), the normed vector space \(R\) (with the norm inherited from \(W\)) is itself a Banach space: by 6.16(b), \(R\) is a complete subset of \(W\), meaning every Cauchy sequence in \(R\) converges to an element of \(R\).

Regard \(T\) as a map \(T : V \to R\). It is linear, bounded (the norm on \(R\) is the norm of \(W\), so \(\|Tf\|_R = \|Tf\|_W \le \|T\| \, \|f\|\)), one-to-one by hypothesis, and onto \(R\) by the definition of \(R\). Because \(V\) and \(R\) are Banach spaces, the Bounded Inverse Theorem (6.83) implies that \(T^{-1} : R \to V\) is a bounded linear map. Thus there exists \(b \in [0, \infty)\) such that

\begin{equation*} \|T^{-1} g\| \le b \|g\| \quad \text{for all } g \in R . \end{equation*}

Let \(c = \max\{b, 1\} \in (0, \infty)\). For \(f \in V\), applying the inequality above to \(g = Tf \in R\) and using \(T^{-1}(Tf) = f\) gives

\begin{equation*} \|f\| \le b \|Tf\| \le c \|Tf\| . \end{equation*}

Hence \(T\) is bounded below, completing the proof.

Problem (6E.13)

Give an example of a Banach space \(V\), a normed vector space \(W\), and a one-to-one bounded linear map \(T\) of \(V\) onto \(W\) such that \(T^{-1}\) is not a bounded linear map of \(W\) onto \(V\).

[This exercise shows that the hypothesis in the Bounded Inverse Theorem (6.83) that \(W\) is a Banach space cannot be relaxed to the hypothesis that \(W\) is a normed vector space.]

Solution

Take the same example as in Exercise 10:

\begin{equation*} V = \ell^1 \text{ with the norm } \|(a_1, a_2, \ldots)\|_1 = \sum_{k=1}^{\infty} |a_k| , \end{equation*}

\begin{equation*} W = \ell^1 \text{ with the norm } \|(a_1, a_2, \ldots)\|_{\infty} = \sup_{k \in \mathbf{Z}^+} |a_k| , \end{equation*}

and let \(T : V \to W\) be the identity map \(Ta = a\).

By the second bullet point of Example 6.38, \(V\) is a Banach space; by the second bullet point of Example 6.39, \(W\) is a normed vector space that is not a Banach space (as it must not be, by 6.83).

\(T\) is linear, one-to-one, and onto \(W\), being the identity map on the common underlying vector space \(\ell^1\). \(T\) is bounded because

\begin{equation*} \|Ta\|_{\infty} = \sup_{k \in \mathbf{Z}^+} |a_k| \le \sum_{k=1}^{\infty} |a_k| = \|a\|_1 \end{equation*}

for all \(a \in \ell^1\); thus \(\|T\| \le 1\).

The inverse map \(T^{-1} : W \to V\) is again the identity map on \(\ell^1\), and it is a linear map of \(W\) onto \(V\). However, \(T^{-1}\) is not bounded. To see this, for \(n \in \mathbf{Z}^+\) let

\begin{equation*} a^{(n)} = \Big( \underbrace{1, \ldots, 1}_{n \text{ terms}}, 0, 0, \ldots \Big) \in \ell^1 . \end{equation*}

Then

\begin{equation*} \|a^{(n)}\|_{\infty} = 1 \quad \text{and} \quad \|T^{-1} a^{(n)}\|_1 = \|a^{(n)}\|_1 = n . \end{equation*}

If there existed \(c \in [0, \infty)\) with \(\|T^{-1} g\|_1 \le c \|g\|_{\infty}\) for all \(g \in W\), then taking \(g = a^{(n)}\) would give \(n \le c\) for every \(n \in \mathbf{Z}^+\), which is impossible. Hence \(T^{-1}\) is an unbounded (equivalently, by 6.48, discontinuous) linear map of \(W\) onto \(V\).

There is no conflict with the Bounded Inverse Theorem (6.83) because \(W\) is not a Banach space: the sequence \(b^{(1)}, b^{(2)}, \ldots\) defined by \(b^{(n)} = (1, \frac{1}{2}, \ldots, \frac{1}{n}, 0, 0, \ldots)\) satisfies \(\|b^{(n)} - b^{(m)}\|_{\infty} = \frac{1}{m+1}\) for \(n > m\), so it is a Cauchy sequence in \(W\); but it has no limit in \(W\), because convergence in the norm \(\|\cdot\|_{\infty}\) implies convergence in each coordinate, which would force the limit to be \((1, \frac{1}{2}, \frac{1}{3}, \ldots)\), a sequence not in \(\ell^1\).

Problem (6E.14)

Show that there exists a normed space \(V\), a Banach space \(W\), and a one-to-one bounded linear map \(T\) of \(V\) onto \(W\) such that \(T^{-1}\) is not a bounded linear map of \(W\) onto \(V\).

[This exercise shows that the hypothesis in the Bounded Inverse Theorem (6.83) that \(V\) is a Banach space cannot be relaxed to the hypothesis that \(V\) is a normed vector space.]

Solution

Take the same construction as in Exercise 11.

Let \(W = \ell^1\) with the norm \(\|(a_1, a_2, \ldots)\|_1 = \sum_{k=1}^{\infty}|a_k|\), which is a Banach space by the second bullet point of Example 6.38. As noted in Exercise 11, \(W\) is infinite-dimensional, so by 6.62 there exists a discontinuous, equivalently (by 6.48) unbounded, linear functional \(\varphi : W \to \mathbf{F}\); thus

\begin{equation*} \text{there is no } c \in (0, \infty) \text{ such that } |\varphi(f)| \le c \|f\|_1 \text{ for all } f \in W . \end{equation*}

Let \(V\) be the vector space \(\ell^1\) equipped with

\begin{equation*} \|f\|_V = \|f\|_1 + |\varphi(f)| . \end{equation*}

As verified in Exercise 11, this is a norm (homogeneity and the triangle inequality follow from those properties for \(\|\cdot\|_1\) and for \(|\varphi(\cdot)|\), using the linearity of \(\varphi\); and \(\|f\|_V = 0\) forces \(\|f\|_1 = 0\), hence \(f = 0\)). So \(V\) is a normed vector space.

Let \(T : V \to W\) be the identity map \(Tf = f\). Then \(T\) is linear, one-to-one, and onto \(W\), because \(V\) and \(W\) have the same underlying vector space \(\ell^1\). Moreover \(T\) is bounded, since

\begin{equation*} \|Tf\|_1 = \|f\|_1 \le \|f\|_1 + |\varphi(f)| = \|f\|_V \end{equation*}

for all \(f \in V\); thus \(\|T\| \le 1\).

The inverse \(T^{-1} : W \to V\) is again the identity map, and it is a linear map of \(W\) onto \(V\). It is not bounded: if there existed \(c \in (0, \infty)\) with \(\|T^{-1}g\|_V \le c\|g\|_1\) for all \(g \in W\), then for every \(f \in \ell^1\) we would have

\begin{equation*} |\varphi(f)| \le \|f\|_1 + |\varphi(f)| = \|f\|_V \le c \|f\|_1 , \end{equation*}

making \(\varphi\) bounded, contrary to the choice of \(\varphi\). Hence \(T^{-1}\) is not a bounded linear map of \(W\) onto \(V\).

Thus \(V\), \(W\), \(T\) have all the required properties. Necessarily \(V\) is not a Banach space: if it were, then the Bounded Inverse Theorem (6.83) applied to the one-to-one bounded linear map \(T\) of the Banach space \(V\) onto the Banach space \(W\) would make \(T^{-1}\) bounded, contradicting what we just proved.

Problem (6E.15)

Prove 6.84.

Solution

Recall the statement of 6.84: if \(V\) and \(W\) are Banach spaces, then \(V \times W\) is a Banach space when given the norm

\begin{equation*} \|(f,g)\| = \max\{\|f\|, \|g\|\} \end{equation*}

for \(f \in V\) and \(g \in W\); moreover, with this norm a sequence \((f_1,g_1), (f_2,g_2), \ldots\) in \(V \times W\) converges to \((f,g)\) if and only if \(\lim_{k \to \infty} f_k = f\) and \(\lim_{k \to \infty} g_k = g\).

Throughout, \(V \times W\) carries the coordinatewise vector space operations

\begin{equation*} (f_1,g_1) + (f_2,g_2) = (f_1+f_2,\; g_1+g_2), \qquad \lambda(f,g) = (\lambda f, \lambda g), \end{equation*}

which make \(V \times W\) into a vector space over \(\mathbf{F}\) (Exercise 10 in Section 6B).

Step 1: the formula defines a norm.

Positive definiteness: for \((f,g) \in V \times W\) we have \(\|(f,g)\| = \max\{\|f\|,\|g\|\} \ge 0\), and \(\|(f,g)\| = 0\) if and only if \(\|f\| = 0\) and \(\|g\| = 0\), which (because \(\|\cdot\|\) is a norm on \(V\) and on \(W\)) happens if and only if \(f = 0\) and \(g = 0\), that is, if and only if \((f,g)\) is the zero element of \(V \times W\).

Homogeneity: if \(\lambda \in \mathbf{F}\) then

\begin{equation*} \|\lambda(f,g)\| = \max\{\|\lambda f\|, \|\lambda g\|\} = \max\{|\lambda|\,\|f\|, |\lambda|\,\|g\|\} = |\lambda| \max\{\|f\|,\|g\|\} = |\lambda|\,\|(f,g)\|, \end{equation*}

where the third equality holds because multiplication by the nonnegative number \(|\lambda|\) preserves the order of real numbers and hence commutes with taking a maximum.

Triangle inequality: suppose \((f_1,g_1), (f_2,g_2) \in V \times W\). Then

\begin{equation*} \|f_1 + f_2\| \le \|f_1\| + \|f_2\| \le \|(f_1,g_1)\| + \|(f_2,g_2)\| \end{equation*}

and similarly

\begin{equation*} \|g_1 + g_2\| \le \|g_1\| + \|g_2\| \le \|(f_1,g_1)\| + \|(f_2,g_2)\|. \end{equation*}

Taking the maximum of the two left sides gives

\begin{equation*} \|(f_1,g_1) + (f_2,g_2)\| = \max\{\|f_1+f_2\|, \|g_1+g_2\|\} \le \|(f_1,g_1)\| + \|(f_2,g_2)\|. \end{equation*}

Thus \(\|\cdot\|\) as defined above is a norm on \(V \times W\), so \(V \times W\) is a normed vector space.

Step 2: convergence in \(V \times W\) is coordinatewise convergence.

Suppose \((f_1,g_1), (f_2,g_2), \ldots\) is a sequence in \(V \times W\) and \((f,g) \in V \times W\). For each \(k \in \mathbf{Z}^+\),

\begin{equation*} \max\{\|f_k - f\|,\; \|g_k - g\|\} = \|(f_k,g_k) - (f,g)\|, \end{equation*}

because \((f_k,g_k) - (f,g) = (f_k - f, g_k - g)\). For nonnegative real numbers \(a\) and \(b\) we have

\begin{equation*} a \le \max\{a,b\} \quad\text{and}\quad b \le \max\{a,b\} \quad\text{and}\quad \max\{a,b\} \le a + b . \end{equation*}

Applying this with \(a = \|f_k - f\|\) and \(b = \|g_k - g\|\) shows that \(\|(f_k,g_k) - (f,g)\| \to 0\) as \(k \to \infty\) if and only if both \(\|f_k - f\| \to 0\) and \(\|g_k - g\| \to 0\). In other words, \(\lim_{k \to \infty} (f_k,g_k) = (f,g)\) in \(V \times W\) if and only if \(\lim_{k \to \infty} f_k = f\) in \(V\) and \(\lim_{k \to \infty} g_k = g\) in \(W\), as claimed.

Step 3: \(V \times W\) is complete.

Suppose \((f_1,g_1), (f_2,g_2), \ldots\) is a Cauchy sequence in \(V \times W\). Let \(\varepsilon > 0\). Then there exists \(n \in \mathbf{Z}^+\) such that

\begin{equation*} \max\{\|f_j - f_k\|, \|g_j - g_k\|\} = \|(f_j,g_j) - (f_k,g_k)\| < \varepsilon \end{equation*}

for all \(j, k \ge n\). In particular \(\|f_j - f_k\| < \varepsilon\) and \(\|g_j - g_k\| < \varepsilon\) for all \(j,k \ge n\). Hence \(f_1, f_2, \ldots\) is a Cauchy sequence in \(V\) and \(g_1, g_2, \ldots\) is a Cauchy sequence in \(W\).

Because \(V\) and \(W\) are Banach spaces, there exist \(f \in V\) and \(g \in W\) with \(\lim_{k \to \infty} f_k = f\) and \(\lim_{k \to \infty} g_k = g\). By Step 2, \(\lim_{k \to \infty} (f_k,g_k) = (f,g)\) in \(V \times W\).

Thus every Cauchy sequence in \(V \times W\) converges to an element of \(V \times W\), which means \(V \times W\) is a Banach space. This completes the proof of 6.84.

Problem (6E.16)

Suppose \(V\) is a Banach space with norm \(\|\cdot\|\) and that \(\varphi : V \to \mathbf{F}\) is a linear functional. Define another norm \(\|\cdot\|_\varphi\) on \(V\) by

\begin{equation*} \|f\|_\varphi = \|f\| + |\varphi(f)|. \end{equation*}

Prove that if \(V\) is a Banach space with the norm \(\|\cdot\|_\varphi\), then \(\varphi\) is a continuous linear functional on \(V\) (with the original norm).

Solution

First we check that \(\|\cdot\|_\varphi\) really is a norm on \(V\), as the exercise asserts. If \(f \in V\) and \(\|f\|_\varphi = 0\), then \(\|f\| + |\varphi(f)| = 0\); both summands are nonnegative, so \(\|f\| = 0\) and hence \(f = 0\). If \(\lambda \in \mathbf{F}\) and \(f \in V\), then the linearity of \(\varphi\) gives

\begin{equation*} \|\lambda f\|_\varphi = \|\lambda f\| + |\varphi(\lambda f)| = |\lambda|\,\|f\| + |\lambda|\,|\varphi(f)| = |\lambda|\,\|f\|_\varphi . \end{equation*}

If \(f, h \in V\), then

\begin{equation*} \|f + h\|_\varphi = \|f+h\| + |\varphi(f) + \varphi(h)| \le \|f\| + \|h\| + |\varphi(f)| + |\varphi(h)| = \|f\|_\varphi + \|h\|_\varphi . \end{equation*}

Thus \(\|\cdot\|_\varphi\) is a norm on \(V\).

Now assume that \(V\) is a Banach space with respect to \(\|\cdot\|_\varphi\) (as well as with respect to \(\|\cdot\|\), which is part of the hypothesis). To distinguish the two normed vector spaces, write \(V_\varphi\) for the vector space \(V\) equipped with the norm \(\|\cdot\|_\varphi\) and write \(V\) for the vector space \(V\) equipped with the original norm \(\|\cdot\|\). By hypothesis both \(V_\varphi\) and \(V\) are Banach spaces.

Let \(I : V_\varphi \to V\) be the identity map, so \(I f = f\) for every \(f \in V\). Then \(I\) is obviously linear, injective, and surjective. Moreover, for every \(f \in V\),

\begin{equation*} \|I f\| = \|f\| \le \|f\| + |\varphi(f)| = \|f\|_\varphi , \end{equation*}

so \(I\) is a bounded linear map with \(\|I\| \le 1\).

Thus \(I\) is a one-to-one bounded linear map from the Banach space \(V_\varphi\) onto the Banach space \(V\). The Bounded Inverse Theorem (6.83) implies that \(I^{-1} : V \to V_\varphi\) is a bounded linear map. Let \(c = \|I^{-1}\| < \infty\). Because \(I^{-1} f = f\) for every \(f \in V\), boundedness of \(I^{-1}\) says precisely that

\begin{equation*} \|f\|_\varphi = \|I^{-1} f\|_\varphi \le c \, \|f\| \end{equation*}

for every \(f \in V\). In other words, the two norms \(\|\cdot\|\) and \(\|\cdot\|_\varphi\) are equivalent.

Now for every \(f \in V\) we have

\begin{equation*} |\varphi(f)| \le \|f\| + |\varphi(f)| = \|f\|_\varphi \le c\,\|f\| . \end{equation*}

Hence \(\varphi\) is a bounded linear functional on \(V\) with respect to the original norm, with \(\|\varphi\| \le c\). By 6.48, boundedness of a linear map is equivalent to its continuity, so \(\varphi\) is a continuous linear functional on \(V\) with the original norm, as desired.

Remark. The converse is also true and easy: if \(\varphi\) is continuous on \(V\) with \(\|\varphi\| = c^{\prime}\), then \(\|f\| \le \|f\|_\varphi \le (1 + c^{\prime})\|f\|\), so the two norms are equivalent and \(V_\varphi\) is complete. Thus \(V_\varphi\) is a Banach space if and only if \(\varphi\) is continuous.

Problem (6E.17)

Suppose \(V\) is a Banach space, \(W\) is a normed vector space, and \(T_1, T_2, \ldots\) is a sequence of bounded linear maps from \(V\) to \(W\) such that \(\lim_{k \to \infty} T_k f\) exists for each \(f \in V\). Define \(T : V \to W\) by

\begin{equation*} T f = \lim_{k \to \infty} T_k f \end{equation*}

for \(f \in V\). Prove that \(T\) is a bounded linear map from \(V\) to \(W\).

[This result states that the pointwise limit of a sequence of bounded linear maps on a Banach space is a bounded linear map.]

Solution

The hypothesis guarantees that \(T f\) is a well-defined element of \(W\) for each \(f \in V\) (limits in a metric space are unique, so the formula defines a genuine function).

First we record two standard facts about limits in a normed vector space \(W\), both of which follow from the triangle inequality.

If \(g_k \to g\) and \(h_k \to h\) in \(W\) and \(\lambda \in \mathbf{F}\), then \(g_k + h_k \to g + h\) and \(\lambda g_k \to \lambda g\). Indeed,

\begin{equation*} \|(g_k + h_k) - (g+h)\| \le \|g_k - g\| + \|h_k - h\| \to 0, \qquad \|\lambda g_k - \lambda g\| = |\lambda|\,\|g_k - g\| \to 0 . \end{equation*}

If \(g_k \to g\) in \(W\), then \(\|g_k\| \to \|g\|\) in \(\mathbf{R}\). Indeed, the triangle inequality gives \(\big|\,\|g_k\| - \|g\|\,\big| \le \|g_k - g\| \to 0\).

Linearity of \(T\). Suppose \(f, h \in V\) and \(\lambda \in \mathbf{F}\). Each \(T_k\) is linear, so \(T_k(f + h) = T_k f + T_k h\) for every \(k\). Since \(T_k f \to T f\) and \(T_k h \to T h\), the first fact above gives \(T_k(f+h) \to Tf + Th\). On the other hand \(T_k(f+h) \to T(f+h)\) by the definition of \(T\). Uniqueness of limits yields

\begin{equation*} T(f+h) = Tf + Th . \end{equation*}

Similarly \(T_k(\lambda f) = \lambda T_k f \to \lambda\, Tf\) and \(T_k(\lambda f) \to T(\lambda f)\), so \(T(\lambda f) = \lambda\, Tf\). Hence \(T\) is a linear map from \(V\) to \(W\).

Boundedness of \(T\). Fix \(f \in V\). The sequence \(T_1 f, T_2 f, \ldots\) converges in \(W\), and every convergent sequence in a metric space is bounded. Explicitly: because \(T_k f \to Tf\), there exists \(n \in \mathbf{Z}^+\) such that \(\|T_k f - Tf\| \le 1\), and hence \(\|T_k f\| \le \|Tf\| + 1\), for all \(k \ge n\). Therefore

\begin{equation*} \sup_{k \in \mathbf{Z}^+} \|T_k f\| \le \max\big(\{\|T_j f\| : 1 \le j < n\} \cup \{\|Tf\| + 1\}\big) < \infty , \end{equation*}

the maximum being over a finite set of real numbers (if \(n = 1\) the first set is empty and the bound is simply \(\|Tf\| + 1\)). Thus, setting \(\mathcal{A} = \{T_k : k \in \mathbf{Z}^+\}\), we have

\begin{equation*} \sup\{\|S f\| : S \in \mathcal{A}\} < \infty \quad \text{for every } f \in V . \end{equation*}

Because \(V\) is a Banach space, \(W\) is a normed vector space, and \(\mathcal{A}\) is a family of bounded linear maps from \(V\) to \(W\) that is pointwise bounded, the Principle of Uniform Boundedness (6.86) implies that

\begin{equation*} M := \sup_{k \in \mathbf{Z}^+} \|T_k\| < \infty . \end{equation*}

Now let \(f \in V\). For each \(k\) we have \(\|T_k f\| \le \|T_k\|\,\|f\| \le M \|f\|\). Since \(T_k f \to Tf\), the second fact recorded above gives \(\|T_k f\| \to \|Tf\|\), and the inequality \(\|T_k f\| \le M\|f\|\) passes to the limit:

\begin{equation*} \|T f\| = \lim_{k \to \infty} \|T_k f\| \le M \|f\| . \end{equation*}

Hence \(T\) is a bounded linear map from \(V\) to \(W\), with \(\|T\| \le M = \sup_k \|T_k\|\).

Remark. The completeness of \(V\) cannot be dropped: on the normed vector space \(V\) of sequences in \(\mathbf{F}\) with only finitely many nonzero terms, with the norm \(\|a\| = \sup_k |a_k|\), the functionals \(T_k a = \sum_{j=1}^{k} a_j\) are bounded and \(\lim_{k \to \infty} T_k a = \sum_{j=1}^{\infty} a_j\) exists for each \(a \in V\) (the sum has only finitely many nonzero terms), but the limit functional is unbounded, because taking \(a\) to be the sequence whose first \(n\) terms equal \(1\) and whose other terms are \(0\) gives \(\|a\| = 1\) and \(Ta = n\).

Problem (6E.18)

Suppose that \(V\) is a normed vector space and \(B\) is a subset of \(V\) such that \(\sup_{f \in B} |\varphi(f)| < \infty\) for every \(\varphi \in V^{\prime}\). Prove that \(\sup_{f \in B} \|f\| < \infty\).

Solution

Note that \(V\) is only assumed to be a normed vector space, so the Principle of Uniform Boundedness (6.86) cannot be applied with \(V\) as the domain. The trick is to apply it with the dual space \(V^{\prime}\) as the domain: by 6.71 and 6.47, \(V^{\prime} = \mathcal{B}(V, \mathbf{F})\) is a Banach space, because the scalar field \(\mathbf{F}\) is complete.

For each \(f \in V\) define \(\widehat{f} : V^{\prime} \to \mathbf{F}\) by

\begin{equation*} \widehat{f}(\varphi) = \varphi(f) \qquad \text{for } \varphi \in V^{\prime}. \end{equation*}

Each \(\widehat{f}\) is linear. Indeed, addition and scalar multiplication in \(V^{\prime}\) are defined pointwise, so for \(\varphi, \psi \in V^{\prime}\) and \(\lambda \in \mathbf{F}\),

\begin{equation*} \widehat{f}(\varphi + \psi) = (\varphi + \psi)(f) = \varphi(f) + \psi(f) = \widehat{f}(\varphi) + \widehat{f}(\psi) \end{equation*}

and

\begin{equation*} \widehat{f}(\lambda \varphi) = (\lambda \varphi)(f) = \lambda\, \varphi(f) = \lambda\, \widehat{f}(\varphi). \end{equation*}

Each \(\widehat{f}\) is bounded, with norm exactly \(\|f\|\). If \(\varphi \in V^{\prime}\), then

\begin{equation*} |\widehat{f}(\varphi)| = |\varphi(f)| \le \|\varphi\|\,\|f\| , \end{equation*}

so \(\widehat{f} \in (V^{\prime})^{\prime}\) with \(\|\widehat{f}\| \le \|f\|\). For the reverse inequality, note first that \(\|\widehat{0}\| = 0 = \|0\|\). If \(f \ne 0\), then 6.72 provides \(\varphi \in V^{\prime}\) with \(\|\varphi\| = 1\) and \(\varphi(f) = \|f\|\); hence

\begin{equation*} \|\widehat{f}\| \ge |\widehat{f}(\varphi)| = |\varphi(f)| = \|f\| . \end{equation*}

Therefore

\begin{equation*} \|\widehat{f}\| = \|f\| \qquad \text{for every } f \in V . \end{equation*}

Now consider the family

\begin{equation*} \mathcal{A} = \{\widehat{f} : f \in B\} \end{equation*}

of bounded linear maps from the Banach space \(V^{\prime}\) to the normed vector space \(\mathbf{F}\). The hypothesis says exactly that \(\mathcal{A}\) is pointwise bounded: for each \(\varphi \in V^{\prime}\),

\begin{equation*} \sup\{|\widehat{f}(\varphi)| : f \in B\} = \sup_{f \in B} |\varphi(f)| < \infty . \end{equation*}

By the Principle of Uniform Boundedness (6.86) applied to the Banach space \(V^{\prime}\) and the family \(\mathcal{A}\),

\begin{equation*} \sup\{\|\widehat{f}\| : f \in B\} < \infty . \end{equation*}

Since \(\|\widehat{f}\| = \|f\|\) for every \(f \in V\), this says

\begin{equation*} \sup_{f \in B} \|f\| < \infty , \end{equation*}

as desired.

Remark. In words: a subset of a normed vector space that is weakly bounded (bounded in the eye of every bounded linear functional) is bounded in norm. The completeness needed to run Baire’s Theorem is supplied for free by the dual space, not by \(V\) itself.

Problem (6E.19)

Suppose \(T : V \to W\) is a linear map from a Banach space \(V\) to a Banach space \(W\) such that

\begin{equation*} \varphi \circ T \in V^{\prime} \quad \text{for all } \varphi \in W^{\prime}. \end{equation*}

Prove that \(T\) is a bounded linear map.

Solution

We use the Closed Graph Theorem (6.85): since \(V\) and \(W\) are Banach spaces, \(T\) is a bounded linear map if and only if

\begin{equation*} \operatorname{graph}(T) = \{(f, Tf) : f \in V\} \end{equation*}

is a closed subspace of \(V \times W\), where \(V \times W\) carries the norm \(\|(f,g)\| = \max\{\|f\|,\|g\|\}\) of 6.84.

\(\operatorname{graph}(T)\) is a subspace. This is immediate from 6.68(a), because \(T\) is assumed to be a linear map.

\(\operatorname{graph}(T)\) is closed. Suppose \((f_1, Tf_1), (f_2, Tf_2), \ldots\) is a sequence in \(\operatorname{graph}(T)\) that converges to some \((f,g) \in V \times W\). By the convergence criterion in 6.84,

\begin{equation*} \lim_{k \to \infty} f_k = f \quad \text{in } V \qquad \text{and} \qquad \lim_{k \to \infty} T f_k = g \quad \text{in } W . \end{equation*}

We must show that \(g = Tf\).

Let \(\varphi \in W^{\prime}\) be arbitrary. On the one hand, \(\varphi\) is a bounded, hence continuous (by 6.48), linear functional on \(W\), so \(T f_k \to g\) implies

\begin{equation*} \lim_{k \to \infty} \varphi(T f_k) = \varphi(g) . \end{equation*}

On the other hand, the hypothesis says \(\varphi \circ T \in V^{\prime}\), so \(\varphi \circ T\) is a bounded, hence continuous, linear functional on \(V\); since \(f_k \to f\) in \(V\), this gives

\begin{equation*} \lim_{k \to \infty} \varphi(T f_k) = \lim_{k \to \infty} (\varphi \circ T)(f_k) = (\varphi \circ T)(f) = \varphi(T f) . \end{equation*}

Uniqueness of limits in \(\mathbf{F}\) now yields

\begin{equation*} \varphi(g) = \varphi(Tf), \qquad \text{that is,} \qquad \varphi(g - Tf) = 0 , \end{equation*}

and this holds for every \(\varphi \in W^{\prime}\).

Suppose, for a contradiction, that \(g - Tf \ne 0\). Then 6.72 provides \(\varphi \in W^{\prime}\) with \(\|\varphi\| = 1\) and \(\varphi(g - Tf) = \|g - Tf\| > 0\), contradicting the previous display. Hence \(g - Tf = 0\), that is, \(g = Tf\).

Therefore \((f,g) = (f, Tf) \in \operatorname{graph}(T)\). Since every convergent sequence in \(\operatorname{graph}(T)\) has its limit in \(\operatorname{graph}(T)\), the set \(\operatorname{graph}(T)\) is closed in \(V \times W\) (by 6.9(e)).

Conclusion. \(\operatorname{graph}(T)\) is a closed subspace of \(V \times W\), so the Closed Graph Theorem (6.85) implies that \(T\) is a bounded linear map from \(V\) to \(W\).

Remark. The Hahn–Banach machinery enters only through 6.72, which supplies enough bounded linear functionals on \(W\) to separate points; this is what converts the pointwise information \(\varphi(g) = \varphi(Tf)\) for all \(\varphi \in W^{\prime}\) into the identity \(g = Tf\).

\(L^p\) Spaces

Exercises 7A

Problem (7A.1)

Suppose \(\mu\) is a measure. Prove that

\begin{equation*} \|f + g\|_\infty \le \|f\|_\infty + \|g\|_\infty \qquad\text{and}\qquad \|\alpha f\|_\infty = |\alpha|\,\|f\|_\infty \end{equation*}

for all \(f, g \in L^\infty(\mu)\) and all \(\alpha \in \mathbf{F}\). Conclude that with the usual operations of addition and scalar multiplication of functions, \(L^\infty(\mu)\) is a vector space.

Solution

Throughout, \((X, \mathcal{S}, \mu)\) is the measure space on which \(\mu\) is defined, and by 7.1

\begin{equation*} \|f\|_\infty = \inf\bigl\{ t > 0 : \mu(\{x \in X : |f(x)| > t\}) = 0 \bigr\}, \end{equation*}

with the convention that the infimum of the empty set is \(\infty\).

Step 1: the infimum is attained, in the sense that \(\mu(\{|f| > \|f\|_\infty\}) = 0\) whenever \(\|f\|_\infty < \infty\).

Write \(c = \|f\|_\infty < \infty\). For each \(n \in \mathbf{Z}^+\), the number \(c + \frac1n\) is strictly larger than the infimum \(c\), so there exists \(t_n > 0\) with \(t_n < c + \frac1n\) and \(\mu(\{x : |f(x)| > t_n\}) = 0\). Since \(t_n < c + \frac1n\) we have

\begin{equation*} \{x : |f(x)| > c + \tfrac1n\} \subseteq \{x : |f(x)| > t_n\}, \end{equation*}

and hence \(\mu(\{x : |f(x)| > c + \frac1n\}) = 0\). Because

\begin{equation*} \{x : |f(x)| > c\} = \bigcup_{n=1}^{\infty} \{x : |f(x)| > c + \tfrac1n\}, \end{equation*}

countable subadditivity of \(\mu\) (2.58) gives \(\mu(\{x : |f(x)| > c\}) = 0\), as claimed. In other words, \(|f(x)| \le \|f\|_\infty\) for almost every \(x \in X\).

Step 2: the triangle inequality.

Suppose \(f, g \in L^\infty(\mu)\), so that \(\|f\|_\infty < \infty\) and \(\|g\|_\infty < \infty\). Let

\begin{equation*} E = \{x \in X : |f(x)| > \|f\|_\infty\} \cup \{x \in X : |g(x)| > \|g\|_\infty\}. \end{equation*}

By Step 1 both sets in this union have measure \(0\), so \(\mu(E) = 0\) by the finite subadditivity that follows from 2.58. If \(x \in X \setminus E\), then

\begin{equation*} |f(x) + g(x)| \le |f(x)| + |g(x)| \le \|f\|_\infty + \|g\|_\infty . \end{equation*}

Hence for every \(t > \|f\|_\infty + \|g\|_\infty\) we have

\begin{equation*} \{x \in X : |f(x) + g(x)| > t\} \subseteq E, \end{equation*}

and therefore \(\mu(\{x : |(f+g)(x)| > t\}) = 0\). Thus every \(t > \|f\|_\infty + \|g\|_\infty\) belongs to the set over which the infimum defining \(\|f + g\|_\infty\) is taken, which forces

\begin{equation*} \|f + g\|_\infty \le \|f\|_\infty + \|g\|_\infty . \end{equation*}

In particular \(\|f+g\|_\infty < \infty\). Also \(f + g\) is \(\mathcal{S}\)-measurable: when \(\mathbf{F} = \mathbf{R}\) this is 2.46, and when \(\mathbf{F} = \mathbf{C}\) it follows from 2.46 applied to \(\operatorname{Re}(f+g) = \operatorname{Re} f + \operatorname{Re} g\) and \(\operatorname{Im}(f+g) = \operatorname{Im} f + \operatorname{Im} g\) together with the definition 6.19 of measurability for complex-valued functions. Hence \(f + g \in L^\infty(\mu)\).

Step 3: homogeneity.

Let \(\alpha \in \mathbf{F}\) and \(f \in L^\infty(\mu)\). If \(\alpha = 0\), then \(\alpha f\) is the function that is identically \(0\), and \(\{x : |0| > t\} = \varnothing\) for every \(t > 0\); hence the infimum defining \(\|\alpha f\|_\infty\) is taken over all of \((0, \infty)\) and equals \(0\). Since \(\|f\|_\infty < \infty\), also \(|\alpha| \|f\|_\infty = 0 \cdot \|f\|_\infty = 0\). So the equality holds in this case.

Now suppose \(\alpha \ne 0\). For every \(t > 0\),

\begin{equation*} \{x \in X : |\alpha f(x)| > t\} = \Bigl\{x \in X : |f(x)| > \tfrac{t}{|\alpha|}\Bigr\}. \end{equation*}

Consequently, writing

\begin{equation*} A_f = \{t > 0 : \mu(\{|f| > t\}) = 0\}, \end{equation*}

the displayed set equality shows that \(t \in A_{\alpha f}\) if and only if \(t/|\alpha| \in A_f\), i.e. \(A_{\alpha f} = |\alpha| \, A_f\). Taking infima and using that multiplication by the positive constant \(|\alpha|\) commutes with taking infima of subsets of \((0,\infty)\), we get

\begin{equation*} \|\alpha f\|_\infty = \inf A_{\alpha f} = |\alpha| \inf A_f = |\alpha| \, \|f\|_\infty . \end{equation*}

(This also covers the case \(A_f = \varnothing\), where both sides are \(\infty\).) In particular \(\|\alpha f\|_\infty < \infty\), and \(\alpha f\) is \(\mathcal{S}\)-measurable (by 2.44 and 2.41, composing \(f\) with the continuous function \(t \mapsto \alpha t\), when \(\mathbf{F} = \mathbf{R}\); for \(\mathbf{F} = \mathbf{C}\) write \(\operatorname{Re}(\alpha f)\) and \(\operatorname{Im}(\alpha f)\) as real linear combinations of \(\operatorname{Re} f\) and \(\operatorname{Im} f\) and use 2.46), so \(\alpha f \in L^\infty(\mu)\).

Step 4: \(L^\infty(\mu)\) is a vector space.

The constant function \(0\) is \(\mathcal{S}\)-measurable and has \(\|0\|_\infty = 0 < \infty\), so \(0 \in L^\infty(\mu)\). By Step 2, \(L^\infty(\mu)\) is closed under addition, and by Step 3 it is closed under scalar multiplication. Thus \(L^\infty(\mu)\) is a subspace of the vector space \(\mathbf{F}^X\) of all functions from \(X\) to \(\mathbf{F}\), and hence is itself a vector space with the usual operations of addition and scalar multiplication of functions.

Problem (7A.2)

Suppose \(a \ge 0\), \(b \ge 0\), and \(1 < p < \infty\). Prove that

\begin{equation*} ab = \frac{a^p}{p} + \frac{b^{p^{\prime}}}{p^{\prime}} \end{equation*}

if and only if \(a^p = b^{p^{\prime}}\) [compare to Young’s inequality (7.8)].

Solution

Recall from 7.6 that \(p^{\prime}\) is determined by \(\frac1p + \frac{1}{p^{\prime}} = 1\); solving gives

\begin{equation*} p^{\prime} = \frac{p}{p-1}, \qquad \frac{p^{\prime}}{p} = \frac{1}{p-1}, \qquad p^{\prime} - 1 = \frac{1}{p-1}. \end{equation*}

First dispose of the degenerate case \(b = 0\). Then the asserted equation reads \(0 = \frac{a^p}{p}\), which holds if and only if \(a = 0\); and the condition \(a^p = b^{p^{\prime}} = 0\) also holds if and only if \(a = 0\). So the two statements are equivalent when \(b = 0\).

Now assume \(b > 0\), and define \(f : [0, \infty) \to \mathbf{R}\) by

\begin{equation*} f(a) = \frac{a^p}{p} + \frac{b^{p^{\prime}}}{p^{\prime}} - ab, \end{equation*}

exactly as in the proof of Young’s inequality (7.8). Thus the equation in question says precisely that \(f(a) = 0\). The function \(f\) is continuous on \([0,\infty)\) and differentiable on \((0, \infty)\), with

\begin{equation*} f^{\prime}(a) = a^{p-1} - b . \end{equation*}

Put \(a_0 = b^{1/(p-1)} > 0\). Since \(p - 1 > 0\), the function \(a \mapsto a^{p-1}\) is strictly increasing on \((0,\infty)\); hence

\begin{equation*} f^{\prime}(a) < 0 \ \text{ for } 0 < a < a_0, \qquad f^{\prime}(a) > 0 \ \text{ for } a > a_0 . \end{equation*}

Therefore \(f\) is strictly decreasing on \([0, a_0]\) and strictly increasing on \([a_0, \infty)\), and so

\begin{equation*} f(a) > f(a_0) \quad \text{for every } a \in [0,\infty) \text{ with } a \ne a_0 . \end{equation*}

Next compute \(f(a_0)\). We have \(a_0^{\,p} = b^{p/(p-1)} = b^{p^{\prime}}\) and \(a_0 b = b^{\frac{1}{p-1} + 1} = b^{\frac{p}{p-1}} = b^{p^{\prime}}\). Hence

\begin{equation*} f(a_0) = \frac{b^{p^{\prime}}}{p} + \frac{b^{p^{\prime}}}{p^{\prime}} - b^{p^{\prime}} = b^{p^{\prime}}\Bigl(\frac1p + \frac{1}{p^{\prime}}\Bigr) - b^{p^{\prime}} = b^{p^{\prime}} - b^{p^{\prime}} = 0 . \end{equation*}

Combining the last two displays: \(f(a) \ge 0\) for all \(a \ge 0\) (which is Young’s inequality 7.8), with

\begin{equation*} f(a) = 0 \iff a = a_0 = b^{1/(p-1)} . \end{equation*}

It remains only to check that, for \(a \ge 0\) and \(b > 0\),

\begin{equation*} a = b^{1/(p-1)} \iff a^p = b^{p^{\prime}} . \end{equation*}

If \(a = b^{1/(p-1)}\), then \(a^p = b^{p/(p-1)} = b^{p^{\prime}}\). Conversely, if \(a^p = b^{p^{\prime}}\), then \(a > 0\) (since \(b^{p^{\prime}} > 0\)) and, raising both sides to the power \(\frac1p\),

\begin{equation*} a = b^{p^{\prime}/p} = b^{1/(p-1)} . \end{equation*}

Putting the pieces together: for \(b > 0\) the equation \(ab = \frac{a^p}{p} + \frac{b^{p^{\prime}}}{p^{\prime}}\) holds if and only if \(f(a) = 0\), if and only if \(a = b^{1/(p-1)}\), if and only if \(a^p = b^{p^{\prime}}\). Together with the case \(b = 0\) handled at the start, this proves the assertion for all \(a, b \ge 0\).

Problem (7A.3)

Suppose \(a_1, \ldots, a_n\) are nonnegative numbers. Prove that

\begin{equation*} (a_1 + \cdots + a_n)^5 \le n^4 \bigl(a_1^{\,5} + \cdots + a_n^{\,5}\bigr). \end{equation*}

Solution

Let \(X = \{1, \ldots, n\}\), let \(\mathcal{S}\) be the \(\sigma\)-algebra of all subsets of \(X\), and let \(\mu\) be counting measure on \(X\). Every function \(X \to \mathbf{R}\) is \(\mathcal{S}\)-measurable, and integration with respect to counting measure is summation, exactly as in 3.6. Indeed, for \(g : X \to [0,\infty)\) the singletons \(\{1\}, \ldots, \{n\}\) are disjoint sets in \(\mathcal{S}\) with \(g = \sum_{k=1}^n g(k) \chi_{\{k\}}\), so 3.7 gives

\begin{equation*} \int g \, d\mu = \sum_{k=1}^{n} g(k) \, \mu(\{k\}) = \sum_{k=1}^{n} g(k). \end{equation*}

Take \(p = 5\). Then by 7.6,

\begin{equation*} p^{\prime} = \frac{p}{p-1} = \frac{5}{4}. \end{equation*}

Define \(f, h : X \to \mathbf{R}\) by \(f(k) = a_k\) and \(h(k) = 1\) for \(k = 1, \ldots, n\). Then

\begin{equation*} \|f h\|_1 = \int |f h| \, d\mu = \sum_{k=1}^{n} a_k , \end{equation*}

\begin{equation*} \|f\|_5 = \Bigl( \int |f|^5 \, d\mu \Bigr)^{1/5} = \Bigl( \sum_{k=1}^{n} a_k^{\,5} \Bigr)^{1/5}, \qquad \|h\|_{5/4} = \Bigl( \sum_{k=1}^{n} 1 \Bigr)^{4/5} = n^{4/5}. \end{equation*}

All three quantities are finite, so Hölder’s inequality (7.9), applied with \(p = 5\) and \(p^{\prime} = 5/4\), gives

\begin{equation*} \sum_{k=1}^{n} a_k \le n^{4/5} \Bigl( \sum_{k=1}^{n} a_k^{\,5} \Bigr)^{1/5}. \end{equation*}

Both sides of this inequality are nonnegative, and \(t \mapsto t^5\) is increasing on \([0,\infty)\), so raising both sides to the fifth power preserves the inequality:

\begin{equation*} \Bigl( \sum_{k=1}^{n} a_k \Bigr)^{5} \le n^{4} \sum_{k=1}^{n} a_k^{\,5}, \end{equation*}

which is the desired inequality.

An alternative proof avoiding Hölder’s inequality: the function \(t \mapsto t^5\) is convex on \([0,\infty)\) because its second derivative \(20t^3\) is nonnegative there, so the finite form of Jensen’s inequality (which follows by induction on \(n\) from the definition of convexity) applied to the equal weights \(\frac1n, \ldots, \frac1n\) gives

\begin{equation*} \Bigl( \frac{a_1 + \cdots + a_n}{n} \Bigr)^{5} \le \frac{a_1^{\,5} + \cdots + a_n^{\,5}}{n} . \end{equation*}

Multiplying both sides by \(n^5\) yields the same conclusion.

Finally, note that the constant \(n^4\) cannot be improved: if \(a_1 = \cdots = a_n = 1\), then the left side equals \(n^5\) and the right side equals \(n^4 \cdot n = n^5\), so equality holds.

Problem (7A.4)

Prove Hölder’s inequality (7.9) in the cases \(p = 1\) and \(p = \infty\).

Solution

Throughout, \((X, \mathcal{S}, \mu)\) is a measure space and \(f, h : X \to \mathbf{F}\) are \(\mathcal{S}\)-measurable; we must show

\begin{equation*} \|f h\|_1 \le \|f\|_p \, \|h\|_{p^{\prime}} . \end{equation*}

By 7.7 we have \(1^{\prime} = \infty\) and \(\infty^{\prime} = 1\), so the two cases to be proved are

\begin{equation*} \|f h\|_1 \le \|f\|_1 \, \|h\|_\infty \qquad\text{and}\qquad \|f h\|_1 \le \|f\|_\infty \, \|h\|_1 . \end{equation*}

Since \(|f h| = |h f|\), the second of these is the first with the roles of \(f\) and \(h\) interchanged. So it suffices to prove the case \(p = 1\), \(p^{\prime} = \infty\); the case \(p = \infty\), \(p^{\prime} = 1\) then follows immediately by that interchange.

Note also that \(f h\) is \(\mathcal{S}\)-measurable (by 2.46 when \(\mathbf{F} = \mathbf{R}\); when \(\mathbf{F} = \mathbf{C}\), expand \(\operatorname{Re}(fh)\) and \(\operatorname{Im}(fh)\) in terms of the real and imaginary parts of \(f\) and \(h\) and apply 2.46 and 6.19), and hence \(|f|\), \(|h|\), \(|f h|\) are \(\mathcal{S}\)-measurable by 6.20 with \(p = 1\). So all the integrals below are defined.

We will use three elementary facts. First, if \(E \in \mathcal{S}\) with \(\mu(E) = 0\) and \(u : X \to [0, \infty]\) is \(\mathcal{S}\)-measurable with \(u(x) = 0\) for all \(x \in X \setminus E\), then \(u \le \infty \cdot \chi_E\) pointwise, so 3.8 and 3.7 give

\begin{equation*} \int u \, d\mu \le \int \infty \cdot \chi_E \, d\mu = \infty \cdot \mu(E) = \infty \cdot 0 = 0, \end{equation*}

hence \(\int u \, d\mu = 0\). Second, for \(c \in [0, \infty)\) and \(\mathcal{S}\)-measurable \(u : X \to [0, \infty]\) we have \(\int c u \, d\mu = c \int u \, d\mu\); this is 3.20, whose first paragraph treats exactly this case of a nonnegative function and a nonnegative real scalar. Third, as shown in Step 1 of the solution to Exercise 1 in this section, if \(\|h\|_\infty < \infty\) then

\begin{equation*} \mu\bigl(\{x \in X : |h(x)| > \|h\|_\infty\}\bigr) = 0 . \end{equation*}

Now suppose \(p = 1\), so \(p^{\prime} = \infty\). We consider several cases.

Case 1: \(\|f\|_1 = 0\). For \(n \in \mathbf{Z}^+\) let \(E_n = \{x \in X : |f(x)| > \frac1n\}\). Then \(\frac1n \chi_{E_n} \le |f|\) pointwise, so 3.8 and 3.7 give \(\frac1n \mu(E_n) \le \int |f| \, d\mu = 0\), and hence \(\mu(E_n) = 0\). Since \(E := \{x : f(x) \ne 0\} = \bigcup_{n=1}^{\infty} E_n\), countable subadditivity (2.58) gives \(\mu(E) = 0\). The function \(|f h|\) vanishes on \(X \setminus E\), so by the first elementary fact above \(\|f h\|_1 = \int |f h| \, d\mu = 0\). Thus \(\|f h\|_1 = 0 \le \|f\|_1 \|h\|_\infty\), whatever value the right side has.

Case 2: \(\|h\|_\infty = 0\). Then by the third fact above, the set \(E = \{x : |h(x)| > 0\}\) has \(\mu(E) = 0\). Again \(|f h|\) vanishes on \(X \setminus E\), so \(\|f h\|_1 = 0\) and the inequality holds.

Case 3: \(\|f\|_1 > 0\) and \(\|h\|_\infty > 0\), and at least one of them equals \(\infty\). Then \(\|f\|_1 \|h\|_\infty = \infty\), and the inequality holds trivially.

Case 4: \(0 < \|f\|_1 < \infty\) and \(0 < \|h\|_\infty < \infty\). Let

\begin{equation*} E = \{x \in X : |h(x)| > \|h\|_\infty\}, \end{equation*}

so \(\mu(E) = 0\) by the third fact above. We claim that

\begin{equation*} |f(x) h(x)| \le \|h\|_\infty |f(x)| + \infty \cdot \chi_E(x) \qquad \text{for all } x \in X . \end{equation*}

Indeed, if \(x \in E\) the right side is \(\infty\), and if \(x \notin E\) then \(|h(x)| \le \|h\|_\infty\), so \(|f(x) h(x)| = |f(x)| \, |h(x)| \le \|h\|_\infty |f(x)|\). Now integrate the displayed inequality: using 3.8 (order preservation), then 3.16 (additivity), then the second elementary fact and 3.7, we obtain

\begin{equation*} \int |f h| \, d\mu \le \int \Bigl( \|h\|_\infty |f| + \infty \cdot \chi_E \Bigr) d\mu = \|h\|_\infty \int |f| \, d\mu + \infty \cdot \mu(E) = \|h\|_\infty \|f\|_1 + 0 . \end{equation*}

That is, \(\|f h\|_1 \le \|f\|_1 \|h\|_\infty\), as desired.

These four cases exhaust all possibilities, so Hölder’s inequality holds when \(p = 1\). As noted at the start, interchanging \(f\) and \(h\) in what we have just proved gives

\begin{equation*} \|f h\|_1 = \|h f\|_1 \le \|h\|_1 \|f\|_\infty = \|f\|_\infty \|h\|_{\infty^{\prime}}, \end{equation*}

which is Hölder’s inequality in the case \(p = \infty\).

Problem (7A.5)

Suppose that \((X, \mathcal{S}, \mu)\) is a measure space, \(1 < p < \infty\), \(f \in L^p(\mu)\), and \(h \in L^{p^{\prime}}(\mu)\). Prove that Hölder’s inequality (7.9) is an equality if and only if there exist nonnegative numbers \(a\) and \(b\), not both \(0\), such that

\begin{equation*} a |f(x)|^p = b |h(x)|^{p^{\prime}} \end{equation*}

for almost every \(x \in X\).

Solution

All the functions integrated below are nonnegative and \(\mathcal{S}\)-measurable: \(f h\) is \(\mathcal{S}\)-measurable (by 2.46 when \(\mathbf{F} = \mathbf{R}\), and via real and imaginary parts together with 6.19 when \(\mathbf{F} = \mathbf{C}\)), and \(|f|^p\), \(|h|^{p^{\prime}}\), \(|f h|\) are then \(\mathcal{S}\)-measurable by 6.20. We also use repeatedly the homogeneity \(\int c u \, d\mu = c \int u \, d\mu\) for \(c \in [0, \infty)\) and \(\mathcal{S}\)-measurable \(u : X \to [0, \infty]\), which is 3.20.

We first record three standard facts, each proved from results in Chapter 3.

(F1) If \(u : X \to [0, \infty]\) is \(\mathcal{S}\)-measurable and \(\int u \, d\mu = 0\), then \(u = 0\) almost everywhere. Indeed, with \(E_n = \{x : u(x) > \frac1n\}\) we have \(\frac1n \chi_{E_n} \le u\), so 3.8 and 3.7 give \(\frac1n \mu(E_n) \le 0\), whence \(\mu(E_n) = 0\); since \(\{x : u(x) > 0\} = \bigcup_{n=1}^\infty E_n\), countable subadditivity (2.58) gives \(\mu(\{u > 0\}) = 0\).

(F2) If \(u, v : X \to [0, \infty]\) are \(\mathcal{S}\)-measurable and \(u = v\) almost everywhere, then \(\int u \, d\mu = \int v \, d\mu\). Indeed, if \(E\) is a set in \(\mathcal{S}\) with \(\mu(E) = 0\) outside of which \(u = v\), then \(u \le v + \infty \cdot \chi_E\) pointwise, so 3.8, 3.16, and 3.7 give \(\int u \, d\mu \le \int v \, d\mu + \infty \cdot \mu(E) = \int v \, d\mu\); the reverse inequality follows by symmetry.

(F3) \(\|f\|_p = 0\) if and only if \(f = 0\) almost everywhere: apply (F1) to \(u = |f|^p\), and conversely apply (F2).

Since \(f \in L^p(\mu)\) and \(h \in L^{p^{\prime}}(\mu)\), both \(\|f\|_p\) and \(\|h\|_{p^{\prime}}\) are finite.

First suppose such \(a, b \ge 0\) exist, not both \(0\), with \(a|f|^p = b|h|^{p^{\prime}}\) almost everywhere. We show \(\|f h\|_1 = \|f\|_p \|h\|_{p^{\prime}}\).

If \(b = 0\), then \(a > 0\) and \(|f|^p = 0\) almost everywhere, so \(f = 0\) almost everywhere; hence \(fh = 0\) almost everywhere, so \(\|f h\|_1 = 0\) by (F2), and \(\|f\|_p = 0\) by (F3), so both sides equal \(0\). The case \(a = 0\) is symmetric (then \(h = 0\) almost everywhere and both sides equal \(0\)). So assume \(a > 0\) and \(b > 0\), and put \(c = b/a > 0\), so that

\begin{equation*} |f(x)|^p = c \, |h(x)|^{p^{\prime}} \quad \text{for almost every } x \in X . \end{equation*}

Integrating and using (F2) gives \(\|f\|_p^{\,p} = c \, \|h\|_{p^{\prime}}^{\,p^{\prime}}\). If \(\|h\|_{p^{\prime}} = 0\) then \(h = 0\) almost everywhere and both sides of Hölder’s inequality equal \(0\), so assume \(\|h\|_{p^{\prime}} > 0\); then also \(\|f\|_p > 0\). Define \(\mathcal{S}\)-measurable functions \(F, H : X \to [0, \infty)\) by

\begin{equation*} F = \frac{|f|}{\|f\|_p}, \qquad H = \frac{|h|}{\|h\|_{p^{\prime}}} . \end{equation*}

Then \(\int F^p \, d\mu = 1\) and \(\int H^{p^{\prime}} \, d\mu = 1\), and for almost every \(x\),

\begin{equation*} F(x)^p = \frac{|f(x)|^p}{\|f\|_p^{\,p}} = \frac{c\,|h(x)|^{p^{\prime}}}{c\,\|h\|_{p^{\prime}}^{\,p^{\prime}}} = H(x)^{p^{\prime}} . \end{equation*}

By Exercise 2 in this section, \(F(x)^p = H(x)^{p^{\prime}}\) implies that Young’s inequality is an equality at \((F(x), H(x))\), i.e.

\begin{equation*} F(x) H(x) = \frac{F(x)^p}{p} + \frac{H(x)^{p^{\prime}}}{p^{\prime}} \quad \text{for almost every } x \in X . \end{equation*}

Integrating this (using (F2) and additivity 3.16, together with homogeneity of the integral of a nonnegative function) gives

\begin{equation*} \int F H \, d\mu = \frac1p \int F^p \, d\mu + \frac{1}{p^{\prime}} \int H^{p^{\prime}} \, d\mu = \frac1p + \frac{1}{p^{\prime}} = 1 . \end{equation*}

Since \(|f h| = \|f\|_p \|h\|_{p^{\prime}} \, F H\), homogeneity gives \(\|f h\|_1 = \|f\|_p \|h\|_{p^{\prime}}\), so Hölder’s inequality is an equality.

Conversely, suppose \(\|f h\|_1 = \|f\|_p \|h\|_{p^{\prime}}\).

If \(\|f\|_p = 0\), take \(a = 1\) and \(b = 0\): by (F3) we have \(f = 0\) almost everywhere, so \(a|f|^p = 0 = b|h|^{p^{\prime}}\) almost everywhere, and \(a, b\) are not both \(0\). Similarly, if \(\|h\|_{p^{\prime}} = 0\), take \(a = 0\) and \(b = 1\).

So assume \(\|f\|_p > 0\) and \(\|h\|_{p^{\prime}} > 0\); as noted, both are finite. Define \(F\) and \(H\) as above, so that \(\int F^p \, d\mu = \int H^{p^{\prime}} \, d\mu = 1\) and, by homogeneity and the assumed equality,

\begin{equation*} \int F H \, d\mu = \frac{\|f h\|_1}{\|f\|_p \|h\|_{p^{\prime}}} = 1 . \end{equation*}

Young’s inequality (7.8) gives, for every \(x \in X\),

\begin{equation*} F(x) H(x) \le \frac{F(x)^p}{p} + \frac{H(x)^{p^{\prime}}}{p^{\prime}} . \end{equation*}

Define \(g : X \to [0, \infty)\) by

\begin{equation*} g = \frac{F^p}{p} + \frac{H^{p^{\prime}}}{p^{\prime}} - F H ; \end{equation*}

this is well defined (all three terms are finite real numbers at each point) and nonnegative, and it is \(\mathcal{S}\)-measurable. Since \(g + FH = \frac{F^p}{p} + \frac{H^{p^{\prime}}}{p^{\prime}}\) with all functions nonnegative, additivity (3.16) and homogeneity give

\begin{equation*} \int g \, d\mu + \int F H \, d\mu = \frac1p \int F^p \, d\mu + \frac{1}{p^{\prime}} \int H^{p^{\prime}} \, d\mu = \frac1p + \frac{1}{p^{\prime}} = 1 . \end{equation*}

Because \(\int F H \, d\mu = 1\) is finite, we may subtract it, obtaining \(\int g \, d\mu = 0\). By (F1), \(g = 0\) almost everywhere, i.e. for almost every \(x \in X\),

\begin{equation*} F(x) H(x) = \frac{F(x)^p}{p} + \frac{H(x)^{p^{\prime}}}{p^{\prime}} . \end{equation*}

By Exercise 2 in this section, this forces \(F(x)^p = H(x)^{p^{\prime}}\) for almost every \(x \in X\), that is,

\begin{equation*} \frac{|f(x)|^p}{\|f\|_p^{\,p}} = \frac{|h(x)|^{p^{\prime}}}{\|h\|_{p^{\prime}}^{\,p^{\prime}}} \quad \text{for almost every } x \in X . \end{equation*}

Multiplying through by \(\|f\|_p^{\,p} \|h\|_{p^{\prime}}^{\,p^{\prime}}\) and setting

\begin{equation*} a = \|h\|_{p^{\prime}}^{\,p^{\prime}} > 0, \qquad b = \|f\|_p^{\,p} > 0, \end{equation*}

we get \(a |f(x)|^p = b |h(x)|^{p^{\prime}}\) for almost every \(x \in X\), with \(a\) and \(b\) nonnegative and not both \(0\). This completes the proof.

Problem (7A.6)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space, \(f \in L^1(\mu)\), and \(h \in L^\infty(\mu)\). Prove that \(\|f h\|_1 = \|f\|_1 \|h\|_\infty\) if and only if

\begin{equation*} |h(x)| = \|h\|_\infty \end{equation*}

for almost every \(x \in X\) such that \(f(x) \ne 0\).

Solution

Write \(M = \|h\|_\infty\) and \(N = \|f\|_1\); both are finite because \(h \in L^\infty(\mu)\) and \(f \in L^1(\mu)\). We use facts (F1) and (F2) established in the solution to Exercise 5 in this section: a nonnegative measurable function with integral \(0\) vanishes almost everywhere, and two nonnegative measurable functions that agree almost everywhere have the same integral.

Reduction. By Step 1 of the solution to Exercise 1 in this section, the set

\begin{equation*} E_0 = \{x \in X : |h(x)| > M\} \end{equation*}

satisfies \(\mu(E_0) = 0\). Let \(\tilde h = h \, \chi_{X \setminus E_0}\). Then \(\tilde h\) is \(\mathcal{S}\)-measurable, \(\tilde h = h\) except on the null set \(E_0\), and \(|\tilde h(x)| \le M\) for every \(x \in X\). Because the essential supremum does not change when a function is modified on a set of measure \(0\) (see the discussion following 7.1), \(\|\tilde h\|_\infty = M\). Also \(|f \tilde h| = |f h|\) except on \(E_0\), so \(\|f \tilde h\|_1 = \|f h\|_1\) by (F2). Finally, the two sets

\begin{equation*} \{x : f(x) \ne 0 \text{ and } |h(x)| \ne M\} \quad\text{and}\quad \{x : f(x) \ne 0 \text{ and } |\tilde h(x)| \ne M\} \end{equation*}

differ only within \(E_0\), so one has measure \(0\) if and only if the other does. Hence both the hypothesis and the conclusion of the assertion are unchanged if \(h\) is replaced by \(\tilde h\), and we may and do assume from now on that

\begin{equation*} |h(x)| \le M \qquad \text{for every } x \in X . \end{equation*}

The key computation. Define \(u : X \to [0, \infty)\) by

\begin{equation*} u(x) = \bigl(M - |h(x)|\bigr) |f(x)| . \end{equation*}

By the normalization just made, \(u \ge 0\), and \(u\) is \(\mathcal{S}\)-measurable. Moreover

\begin{equation*} u + |f h| = M |f| \end{equation*}

pointwise on \(X\), and all three of these functions are nonnegative and \(\mathcal{S}\)-measurable (for \(u\) this uses 2.46 and 6.20; for \(|f h|\) it uses 6.20 applied to \(f h\)). Hence additivity of integration (3.16), together with homogeneity of the integral (3.20, applied with the constant \(M \in [0, \infty)\)), gives

\begin{equation*} \int u \, d\mu + \|f h\|_1 = M \|f\|_1 = M N . \end{equation*}

Since \(\|f h\|_1 \le \|f\|_1 \|h\|_\infty = M N < \infty\) by Hölder’s inequality (7.9) in the case \(p = 1\) (proved in Exercise 4 of this section), the term \(\|f h\|_1\) is finite and may be subtracted, yielding

\begin{equation*} \int u \, d\mu = M N - \|f h\|_1 = \|f\|_1 \|h\|_\infty - \|f h\|_1 . \end{equation*}

Conclusion. From the last display,

\begin{equation*} \|f h\|_1 = \|f\|_1 \|h\|_\infty \iff \int u \, d\mu = 0 . \end{equation*}

By (F1), \(\int u \, d\mu = 0\) if and only if \(u = 0\) almost everywhere. Now observe that, because \(|h| \le M\) everywhere,

\begin{equation*} \{x \in X : u(x) > 0\} = \{x \in X : f(x) \ne 0 \text{ and } |h(x)| < M\} = \{x \in X : f(x) \ne 0 \text{ and } |h(x)| \ne M\} . \end{equation*}

Therefore \(u = 0\) almost everywhere if and only if

\begin{equation*} \mu\bigl(\{x \in X : f(x) \ne 0 \text{ and } |h(x)| \ne \|h\|_\infty\}\bigr) = 0, \end{equation*}

which is exactly the statement that \(|h(x)| = \|h\|_\infty\) for almost every \(x \in X\) such that \(f(x) \ne 0\).

Combining the two displayed equivalences completes the proof: \(\|f h\|_1 = \|f\|_1 \|h\|_\infty\) if and only if \(|h(x)| = \|h\|_\infty\) for almost every \(x \in X\) with \(f(x) \ne 0\). (Undoing the reduction, the same statement holds for the original \(h\), as explained above.)

Problem (7A.7)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(f, h : X \to \mathbf{F}\) are \(\mathcal{S}\)-measurable. Prove that

\begin{equation*} \|f h\|_r \le \|f\|_p \, \|h\|_q \end{equation*}

for all positive numbers \(p, q, r\) such that \(\frac1p + \frac1q = \frac1r\).

Solution

Let \(p, q, r\) be positive (finite) numbers with \(\frac1p + \frac1q = \frac1r\). The idea is to reduce to Hölder’s inequality (7.9) by applying it to the functions \(|f|^r\) and \(|h|^r\) with a suitable exponent.

Choice of exponents. Set

\begin{equation*} P = \frac{p}{r}, \qquad Q = \frac{q}{r}. \end{equation*}

Since \(\frac1q > 0\), the hypothesis gives \(\frac1r = \frac1p + \frac1q > \frac1p\), hence \(r < p\) and therefore \(P > 1\); also \(P < \infty\) because \(r > 0\) and \(p < \infty\). Symmetrically \(1 < Q < \infty\). Furthermore

\begin{equation*} \frac1P + \frac1Q = \frac{r}{p} + \frac{r}{q} = r\Bigl( \frac1p + \frac1q \Bigr) = r \cdot \frac1r = 1 . \end{equation*}

By the definition of the dual exponent (7.6), this says exactly that

\begin{equation*} Q = P^{\prime} . \end{equation*}

Applying Hölder’s inequality. The functions \(|f|^r\) and \(|h|^r\) map \(X\) into \([0, \infty)\) and are \(\mathcal{S}\)-measurable by 6.20; since \([0,\infty) \subseteq \mathbf{F}\), they are legitimate inputs to 7.9. Hölder’s inequality (7.9), applied with exponent \(P \in (1, \infty)\) to the pair \(|f|^r\), \(|h|^r\), gives

\begin{equation*} \bigl\| \, |f|^r \, |h|^r \, \bigr\|_1 \le \bigl\| \, |f|^r \, \bigr\|_P \; \bigl\| \, |h|^r \, \bigr\|_{P^{\prime}} = \bigl\| \, |f|^r \, \bigr\|_P \; \bigl\| \, |h|^r \, \bigr\|_{Q} . \end{equation*}

Identifying the three terms. First,

\begin{equation*} \bigl\| \, |f|^r |h|^r \, \bigr\|_1 = \int |f h|^r \, d\mu = \|f h\|_r^{\,r} , \end{equation*}

directly from Definition 7.1. Second, using \(rP = p\),

\begin{equation*} \bigl\| \, |f|^r \, \bigr\|_P = \Bigl( \int |f|^{rP} \, d\mu \Bigr)^{1/P} = \Bigl( \int |f|^{p} \, d\mu \Bigr)^{r/p} = \Bigl[ \Bigl( \int |f|^{p} \, d\mu \Bigr)^{1/p} \Bigr]^{r} = \|f\|_p^{\,r} . \end{equation*}

Third, in the same way, using \(rQ = q\),

\begin{equation*} \bigl\| \, |h|^r \, \bigr\|_{Q} = \|h\|_q^{\,r} . \end{equation*}

(These identities hold in \([0, \infty]\), with the usual conventions for powers of \(\infty\).)

Conclusion. Substituting the three identities into the Hölder estimate gives

\begin{equation*} \|f h\|_r^{\,r} \le \|f\|_p^{\,r} \, \|h\|_q^{\,r} = \bigl( \|f\|_p \, \|h\|_q \bigr)^{r} , \end{equation*}

where the last equality uses the convention \(0 \cdot \infty = 0\) if one factor is \(0\) and the other is \(\infty\) (the same convention under which Hölder’s inequality 7.9 is stated). Since \(t \mapsto t^{1/r}\) is increasing on \([0, \infty]\), raising both sides to the power \(\frac1r\) preserves the inequality and yields

\begin{equation*} \|f h\|_r \le \|f\|_p \, \|h\|_q , \end{equation*}

as desired.

Remark. If one also allows the value \(\infty\) for \(p\) or \(q\) (with the convention \(\frac{1}{\infty} = 0\)), the result still holds. Suppose \(p = \infty\); then the hypothesis forces \(q = r\), and we claim \(\|f h\|_r \le \|f\|_\infty \|h\|_r\). If \(\|f\|_\infty = \infty\) this is trivial, so assume \(\|f\|_\infty < \infty\) and let \(E = \{x \in X : |f(x)| > \|f\|_\infty\}\), a set with \(\mu(E) = 0\) by Step 1 of the solution to Exercise 1 in this section. Then

\begin{equation*} |f(x) h(x)|^r \le \|f\|_\infty^{\,r} |h(x)|^r + \infty \cdot \chi_E(x) \qquad \text{for all } x \in X, \end{equation*}

since the right side equals \(\infty\) on \(E\) and the inequality is clear off \(E\). Integrating and using 3.8, then 3.16, then 3.20 and 3.7, gives

\begin{equation*} \|f h\|_r^{\,r} \le \|f\|_\infty^{\,r} \|h\|_r^{\,r} + \infty \cdot \mu(E) = \|f\|_\infty^{\,r} \|h\|_r^{\,r}, \end{equation*}

and taking \(r\)-th roots yields the claim. The case \(q = \infty\) is symmetric.

Problem (7A.8)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(n \in \mathbb{Z}^{+}\). Prove that

\begin{equation*} \|f_1 f_2 \cdots f_n\|_1 \le \|f_1\|_{p_1} \|f_2\|_{p_2} \cdots \|f_n\|_{p_n} \end{equation*}

for all positive numbers \(p_1, \ldots, p_n\) such that \(\frac{1}{p_1} + \frac{1}{p_2} + \cdots + \frac{1}{p_n} = 1\) and all \(\mathcal{S}\)-measurable functions \(f_1, f_2, \ldots, f_n : X \to \mathbf{F}\).

Solution

Throughout we use the convention \(0 \cdot \infty = 0\) for the right-hand side, and we allow the exponents to lie in \((0, \infty]\) (with \(1/\infty = 0\)), which includes the stated case. Note first that if \(\frac{1}{p_1} + \cdots + \frac{1}{p_n} = 1\) with each \(p_k > 0\), then \(\frac{1}{p_k} \le 1\) for each \(k\), so each \(p_k \ge 1\) and all the norms appearing above are genuine \(p\)-norms with \(p \ge 1\).

Step 1: the two-factor inequality (this is Exercise 7 in this section). We claim that if \(p, q, r \in (0, \infty]\) satisfy \(\frac{1}{p} + \frac{1}{q} = \frac{1}{r}\), then

\begin{equation*} \|f h\|_r \le \|f\|_p \|h\|_q \end{equation*}

for all \(\mathcal{S}\)-measurable \(f, h : X \to \mathbf{F}\).

We first record an observation used repeatedly: for \(s \in (0, \infty]\) and \(\mathcal{S}\)-measurable \(g\), we have \(\|g\|_s = 0\) if and only if \(g = 0\) almost everywhere. If \(s < \infty\) this is because \(\|g\|_s = 0\) means \(\int |g|^s \, d\mu = 0\), which by 3.43 (applied to the nonnegative function \(|g|^s\)) forces \(|g|^s = 0\) almost everywhere. If \(s = \infty\), then \(\|g\|_\infty = 0\) means that for every \(t > 0\) the set \(\{x : |g(x)| > t\}\) has measure \(0\) (by the definition of the infimum in 7.1 and monotonicity of \(\mu\)), so \(\{x : g(x) \ne 0\} = \bigcup_{m=1}^{\infty} \{x : |g(x)| > \frac{1}{m}\}\) has measure \(0\) by countable subadditivity (2.58). The converse direction is clear in both cases.

Now dispose of the degenerate cases of Step 1, using the convention \(0 \cdot \infty = 0\). If \(\|f\|_p = 0\) or \(\|h\|_q = 0\), then by the previous paragraph \(f = 0\) almost everywhere or \(h = 0\) almost everywhere, so \(fh = 0\) almost everywhere, whence \(\|fh\|_r = 0\) and the asserted inequality holds with both sides equal to \(0\). If neither of \(\|f\|_p, \|h\|_q\) is \(0\) and at least one of them equals \(\infty\), then the right-hand side is \(\infty\) and the inequality is trivial. So we may assume \(0 < \|f\|_p < \infty\) and \(0 < \|h\|_q < \infty\).

First suppose \(q = \infty\); then \(\frac{1}{r} = \frac{1}{p}\), so \(r = p\). Let \(t > \|h\|_{\infty}\). Since by 7.1 the number \(\|h\|_{\infty}\) is the infimum of \(\{s > 0 : \mu(\{x : |h(x)| > s\}) = 0\}\) and \(\|h\|_\infty < t\), there exists \(s\) in that set with \(s < t\); because \(\{x : |h(x)| > t\} \subseteq \{x : |h(x)| > s\}\), this gives \(\mu(\{x : |h(x)| > t\}) = 0\). Applying this with \(t = \|h\|_\infty + \frac{1}{m}\) for \(m \in \mathbb{Z}^{+}\) (legitimate since \(\|h\|_\infty < \infty\) here) and taking the union of the resulting sets of measure \(0\) shows, by countable subadditivity (2.58), that

\begin{equation*} \mu\bigl(\{x \in X : |h(x)| > \|h\|_\infty\}\bigr) = 0 . \end{equation*}

Thus \(|f(x) h(x)| \le \|h\|_\infty |f(x)|\) for almost every \(x \in X\), and therefore \(\|fh\|_p \le \|h\|_\infty \|f\|_p\), which is the claim (this argument also covers \(p = \infty\), in which case \(r = \infty\) too). The case \(p = \infty\) with \(q < \infty\) follows by symmetry.

Now suppose \(p < \infty\) and \(q < \infty\); then \(r < \infty\) as well. Put \(P = \frac{p}{r}\). Then

\begin{equation*} \frac{1}{P} = \frac{r}{p} = r\Bigl(\frac{1}{r} - \frac{1}{q}\Bigr) = 1 - \frac{r}{q} \in (0, 1), \end{equation*}

so \(1 < P < \infty\), and the equation above shows that the dual exponent of \(P\) is \(P^{\prime} = \frac{q}{r}\). Apply Hölder’s inequality (7.9), with exponent \(P\), to the \(\mathcal{S}\)-measurable functions \(|f|^r\) and \(|h|^r\):

\begin{equation*} \int |f|^r |h|^r \, d\mu \le \Bigl(\int |f|^{rP} d\mu\Bigr)^{1/P} \Bigl(\int |h|^{rP^{\prime}} d\mu\Bigr)^{1/P^{\prime}} = \Bigl(\int |f|^{p} d\mu\Bigr)^{r/p} \Bigl(\int |h|^{q} d\mu\Bigr)^{r/q} . \end{equation*}

In other words \(\|fh\|_r^{\,r} \le \|f\|_p^{\,r} \|h\|_q^{\,r}\); raising both sides to the power \(\frac{1}{r}\) gives \(\|fh\|_r \le \|f\|_p \|h\|_q\), completing Step 1.

Step 2: induction. We prove the following statement by induction on \(n\): if \(n \in \mathbb{Z}^{+}\), if \(p_1, \ldots, p_n, r \in (0, \infty]\) satisfy \(\frac{1}{p_1} + \cdots + \frac{1}{p_n} = \frac{1}{r}\), and if \(f_1, \ldots, f_n : X \to \mathbf{F}\) are \(\mathcal{S}\)-measurable, then

\begin{equation*} \|f_1 f_2 \cdots f_n\|_r \le \|f_1\|_{p_1} \|f_2\|_{p_2} \cdots \|f_n\|_{p_n} . \end{equation*}

If \(n = 1\) then \(r = p_1\) and both sides are equal. Suppose now that \(n \in \mathbb{Z}^{+}\) and the statement holds for \(n\); suppose \(p_1, \ldots, p_{n+1}, r \in (0, \infty]\) with \(\sum_{k=1}^{n+1} \frac{1}{p_k} = \frac{1}{r}\). Define \(s \in (0, \infty]\) by

\begin{equation*} \frac{1}{s} = \sum_{k=1}^{n} \frac{1}{p_k} \end{equation*}

(so \(s = \infty\) if that sum is \(0\)). Then \(\frac{1}{s} + \frac{1}{p_{n+1}} = \frac{1}{r}\). By the induction hypothesis,

\begin{equation*} \|f_1 \cdots f_n\|_s \le \|f_1\|_{p_1} \cdots \|f_n\|_{p_n}, \end{equation*}

and by Step 1 applied to the measurable functions \(f_1 \cdots f_n\) and \(f_{n+1}\),

\begin{equation*} \|f_1 \cdots f_n f_{n+1}\|_r \le \|f_1 \cdots f_n\|_s \|f_{n+1}\|_{p_{n+1}} \le \|f_1\|_{p_1} \cdots \|f_n\|_{p_n} \|f_{n+1}\|_{p_{n+1}} . \end{equation*}

(If some \(\|f_k\|_{p_k} = 0\), then \(f_k = 0\) almost everywhere by the observation at the start of Step 1, so \(f_1 \cdots f_{n+1} = 0\) almost everywhere and both sides equal \(0\) under the convention \(0 \cdot \infty = 0\); if no factor vanishes and some factor is \(\infty\), the right side is \(\infty\) and the inequality is trivial. So the multiplications above are unambiguous.) This completes the induction.

Conclusion. Given positive numbers \(p_1, \ldots, p_n\) with \(\frac{1}{p_1} + \cdots + \frac{1}{p_n} = 1\), apply Step 2 with \(r = 1\) to obtain

\begin{equation*} \|f_1 f_2 \cdots f_n\|_1 \le \|f_1\|_{p_1} \|f_2\|_{p_2} \cdots \|f_n\|_{p_n}, \end{equation*}

as desired. Note that the case \(n = 2\) is exactly Hölder’s inequality (7.9), since then \(p_2 = p_1^{\prime}\).

Problem (7A.9)

Show that the formula in 7.12 holds for \(p = \infty\) if \(\mu\) is a \(\sigma\)-finite measure.

Solution

We must show: if \((X, \mathcal{S}, \mu)\) is a \(\sigma\)-finite measure space and \(f \in \mathcal{L}^{\infty}(\mu)\), then

\begin{equation*} \|f\|_{\infty} = \sup\Bigl\{ \int f h \, d\mu : h \in \mathcal{L}^{1}(\mu) \text{ and } \|h\|_1 \le 1 \Bigr\}, \end{equation*}

which is the formula in 7.12 with \(p = \infty\) and \(p^{\prime} = \infty^{\prime} = 1\) (see 7.6 and 7.7). Example 7.13 shows that the \(\sigma\)-finiteness hypothesis cannot be dropped. As in 7.12, the expression \(\int f h \, d\mu\) is read as a real number when \(\mathbf{F} = \mathbb{R}\); when \(\mathbf{F} = \mathbb{C}\) it is bounded above via \(\bigl| \int f h \, d\mu \bigr| \le \int |f h| \, d\mu\), and the functions \(h\) constructed below make \(f h\) nonnegative, so the supremum is unchanged under either reading. The argument below therefore covers both choices of \(\mathbf{F}\).

We first show that the supremum is at most \(\|f\|_{\infty}\). Suppose \(h \in \mathcal{L}^1(\mu)\) with \(\|h\|_1 \le 1\). Hölder’s inequality (7.9) with \(p = \infty\) and \(p^{\prime} = 1\) gives

\begin{equation*} \int f h \, d\mu \le \int |f h| \, d\mu = \|f h\|_1 \le \|f\|_{\infty} \|h\|_1 \le \|f\|_{\infty} . \end{equation*}

Hence the supremum on the right side is at most \(\|f\|_\infty\). Note also that \(h = 0\) is admissible, so the supremum is at least \(0\); in particular, if \(\|f\|_\infty = 0\) then both sides equal \(0\) and we are done. Assume from now on that \(\|f\|_{\infty} > 0\).

We now show that the supremum is at least \(\|f\|_{\infty}\). Let \(t \in (0, \|f\|_{\infty})\) and set

\begin{equation*} E = \{x \in X : |f(x)| > t\}, \end{equation*}

which is in \(\mathcal{S}\) because \(f\) is \(\mathcal{S}\)-measurable. Since \(t < \|f\|_\infty\) and \(\|f\|_\infty\) is the infimum of the set \(\{s > 0 : \mu(\{x : |f(x)| > s\}) = 0\}\) (definition 7.1), the number \(t\) does not belong to that set; thus \(\mu(E) > 0\).

Because \(\mu\) is \(\sigma\)-finite, there exist \(X_1, X_2, \ldots \in \mathcal{S}\) with \(X = \bigcup_{j=1}^{\infty} X_j\) and \(\mu(X_j) < \infty\) for each \(j\). Then \(E = \bigcup_{j=1}^{\infty} (E \cap X_j)\), so if \(\mu(E \cap X_j) = 0\) for every \(j\), countable subadditivity of \(\mu\) (2.58) would give \(\mu(E) = 0\), a contradiction. Hence there exists \(j\) with \(\mu(E \cap X_j) > 0\). Put

\begin{equation*} A = E \cap X_j, \qquad \text{so } 0 < \mu(A) < \infty . \end{equation*}

Define \(h : X \to \mathbf{F}\) by

\begin{equation*} h(x) = \begin{cases} \dfrac{\overline{f(x)}}{\mu(A) \, |f(x)|} & \text{if } x \in A, \\[2mm] 0 & \text{otherwise} \end{cases} \end{equation*}

(when \(\mathbf{F} = \mathbb{R}\) this is \(\frac{1}{\mu(A)}\) times the sign of \(f\) on \(A\)). The definition makes sense because \(|f(x)| > t > 0\) for \(x \in A\), and \(h\) is \(\mathcal{S}\)-measurable. Moreover \(|h| = \frac{1}{\mu(A)} \chi_A\), so

\begin{equation*} \|h\|_1 = \int |h| \, d\mu = \frac{\mu(A)}{\mu(A)} = 1 , \end{equation*}

using 3.15; in particular \(h \in \mathcal{L}^1(\mu)\) and \(h\) is admissible in the supremum. Finally, \(f(x) h(x) = \frac{|f(x)|}{\mu(A)}\) for \(x \in A\) and \(f(x)h(x) = 0\) otherwise, so

\begin{equation*} \int f h \, d\mu = \frac{1}{\mu(A)} \int \chi_A |f| \, d\mu \ge \frac{1}{\mu(A)} \int t \chi_A \, d\mu = \frac{t \, \mu(A)}{\mu(A)} = t , \end{equation*}

where the inequality holds because \(t \chi_A \le \chi_A |f|\) pointwise on \(X\) (as \(|f| > t\) on \(A\)) and integration of nonnegative functions is order preserving (3.8), and the last two equalities use 3.15 together with \(0 < \mu(A) < \infty\).

Thus the supremum is at least \(t\) for every \(t \in (0, \|f\|_{\infty})\), and therefore the supremum is at least \(\|f\|_{\infty}\). Combining the two inequalities gives the desired equality.

Note where \(\sigma\)-finiteness was used: it produced a set \(A\) of finite positive measure inside \(E\), which is what makes the normalized function \(h\) lie in \(\mathcal{L}^1(\mu)\). In Example 7.13 no such set exists, and indeed the formula fails there.

Problem (7A.10)

Suppose \(0 < p < q \le \infty\).

(a)
Prove that \(\ell^p \subseteq \ell^q\).
(b)
Prove that \(\|(a_1, a_2, \ldots)\|_p \ge \|(a_1, a_2, \ldots)\|_q\) for every sequence \(a_1, a_2, \ldots\) of elements of \(\mathbf{F}\).
Solution

Throughout, \(\mu\) is counting measure on \(\mathbb{Z}^{+}\), so by 7.2 we have \(\|a\|_p = \bigl(\sum_{k=1}^{\infty} |a_k|^p\bigr)^{1/p}\) for \(0 < p < \infty\) and \(\|a\|_{\infty} = \sup_{k \in \mathbb{Z}^{+}} |a_k|\), and \(\ell^p = \{a : \|a\|_p < \infty\}\) (7.4).

(a) Suppose \(a = (a_1, a_2, \ldots) \in \ell^p\), so that \(\sum_{k=1}^{\infty} |a_k|^p < \infty\). Since the terms of a convergent series tend to \(0\), we have \(|a_k|^p \to 0\) and hence \(|a_k| \to 0\); in particular the sequence \((|a_k|)\) is bounded. Let

\begin{equation*} M = \sup_{k \in \mathbb{Z}^{+}} |a_k| < \infty . \end{equation*}

If \(q = \infty\) this already gives \(\|a\|_{\infty} = M < \infty\), so \(a \in \ell^{\infty}\).

Now suppose \(q < \infty\). Then for every \(k\),

\begin{equation*} |a_k|^q = |a_k|^p \, |a_k|^{q-p} \le M^{q-p} |a_k|^p , \end{equation*}

because \(q - p > 0\) and \(|a_k| \le M\). Summing over \(k\) gives

\begin{equation*} \sum_{k=1}^{\infty} |a_k|^q \le M^{q-p} \sum_{k=1}^{\infty} |a_k|^p < \infty , \end{equation*}

so \(a \in \ell^q\). In both cases \(a \in \ell^q\), proving \(\ell^p \subseteq \ell^q\).

(b) Let \(a = (a_1, a_2, \ldots)\) be a sequence in \(\mathbf{F}\). If \(\|a\|_p = \infty\) the inequality is trivial, so assume \(\|a\|_p < \infty\), that is, \(a \in \ell^p\). If \(\|a\|_p = 0\) then \(a_k = 0\) for every \(k\), so \(\|a\|_q = 0\) as well and the inequality holds. Hence assume \(0 < \|a\|_p < \infty\) and set

\begin{equation*} b = (b_1, b_2, \ldots) = \frac{1}{\|a\|_p} \, a , \end{equation*}

so that, by the homogeneity of the \(p\)-norm (7.5(b)), \(\|b\|_p = 1\); that is, \(\sum_{k=1}^{\infty} |b_k|^p = 1\). In particular \(|b_k|^p \le 1\), so

\begin{equation*} |b_k| \le 1 \quad \text{for every } k \in \mathbb{Z}^{+} . \end{equation*}

If \(q = \infty\), this says exactly \(\|b\|_{\infty} = \sup_k |b_k| \le 1\). If \(q < \infty\), then since \(0 \le |b_k| \le 1\) and \(q > p\) we have \(|b_k|^q \le |b_k|^p\) for every \(k\), and therefore

\begin{equation*} \|b\|_q^{\,q} = \sum_{k=1}^{\infty} |b_k|^q \le \sum_{k=1}^{\infty} |b_k|^p = 1 , \end{equation*}

which gives \(\|b\|_q \le 1\). So in all cases \(\|b\|_q \le 1\).

Finally, homogeneity of the \(q\)-norm (7.5(b) if \(q < \infty\), Exercise 1 in this section if \(q = \infty\)) gives

\begin{equation*} \|a\|_q = \bigl\| \, \|a\|_p \, b \, \bigr\|_q = \|a\|_p \, \|b\|_q \le \|a\|_p , \end{equation*}

as desired.

Note that (b) gives a second proof of (a): if \(a \in \ell^p\) then \(\|a\|_q \le \|a\|_p < \infty\), so \(a \in \ell^q\). The direction of the containment is opposite to that in 7.10; the reason is that counting measure on \(\mathbb{Z}^{+}\) is not a finite measure, but it has no sets of small positive measure, so the mass of a sequence cannot spread out thinly.

Problem (7A.11)

Show that

\begin{equation*} \bigcap_{p > 1} \ell^p \ne \ell^1 . \end{equation*}

Solution

By Exercise 10(a) in this section, \(\ell^1 \subseteq \ell^p\) for every \(p > 1\) (including \(p = \infty\)), and hence

\begin{equation*} \ell^1 \subseteq \bigcap_{p > 1} \ell^p . \end{equation*}

So it suffices to exhibit a sequence that lies in \(\ell^p\) for every \(p > 1\) but not in \(\ell^1\). Let

\begin{equation*} a = (a_1, a_2, \ldots), \qquad a_k = \frac{1}{k} . \end{equation*}

First, \(a \notin \ell^1\). Indeed, for each \(m \in \{0, 1, 2, \ldots\}\),

\begin{equation*} \sum_{k = 2^m}^{2^{m+1} - 1} \frac{1}{k} \ge 2^m \cdot \frac{1}{2^{m+1}} = \frac{1}{2} , \end{equation*}

because the block contains \(2^m\) terms, each at least \(\frac{1}{2^{m+1}}\). Summing these blocks for \(m = 0, 1, \ldots, M-1\) gives

\begin{equation*} \sum_{k=1}^{2^M - 1} \frac{1}{k} \ge \frac{M}{2} , \end{equation*}

so the partial sums of \(\sum_{k=1}^{\infty} \frac{1}{k}\) are unbounded. Thus \(\|a\|_1 = \infty\) and \(a \notin \ell^1\).

Second, \(a \in \ell^p\) for every \(p > 1\). If \(p = \infty\) this is clear, since \(\sup_k \frac{1}{k} = 1 < \infty\). Now suppose \(1 < p < \infty\). For each \(m \in \{0, 1, 2, \ldots\}\) the block \(\{2^m, \ldots, 2^{m+1}-1\}\) contains \(2^m\) terms, each satisfying \(k^{-p} \le (2^m)^{-p}\); hence

\begin{equation*} \sum_{k = 2^m}^{2^{m+1} - 1} \frac{1}{k^p} \le 2^m \cdot 2^{-mp} = \bigl(2^{1-p}\bigr)^{m} . \end{equation*}

Since \(p > 1\), we have \(0 < 2^{1-p} < 1\), so summing the geometric series gives

\begin{equation*} \sum_{k=1}^{\infty} \frac{1}{k^p} \le \sum_{m=0}^{\infty} \bigl(2^{1-p}\bigr)^{m} = \frac{1}{1 - 2^{1-p}} < \infty \end{equation*}

(the partial sums of the nonnegative series \(\sum k^{-p}\) are bounded by this number, hence the series converges). Thus \(\|a\|_p < \infty\), so \(a \in \ell^p\).

Therefore \(a \in \bigcap_{p>1} \ell^p\) but \(a \notin \ell^1\), which shows that the containment \(\ell^1 \subseteq \bigcap_{p>1} \ell^p\) is strict. In particular \(\bigcap_{p>1} \ell^p \ne \ell^1\), as desired.

Problem (7A.12)

Show that

\begin{equation*} \bigcap_{p < \infty} \mathcal{L}^p([0,1]) \ne \mathcal{L}^{\infty}([0,1]) . \end{equation*}

Solution

Here \(\lambda\) denotes Lebesgue measure on \([0,1]\), and \(\mathcal{L}^p([0,1]) = \mathcal{L}^p(\lambda_{[0,1]})\) as in 7.11; note \(\lambda([0,1]) = 1\).

First we check the containment \(\mathcal{L}^{\infty}([0,1]) \subseteq \bigcap_{p<\infty} \mathcal{L}^p([0,1])\). Suppose \(f \in \mathcal{L}^{\infty}([0,1])\) and \(0 < p < \infty\); note \(\|f\|_\infty < \infty\). As shown at the start of the solution to Exercise 8 (apply the argument there with \(t = \|f\|_\infty + \frac{1}{m}\) for \(m \in \mathbb{Z}^{+}\) and use countable subadditivity 2.58), the set

\begin{equation*} N = \{x \in [0,1] : |f(x)| > \|f\|_{\infty}\} \end{equation*}

satisfies \(\lambda(N) = 0\). Pointwise on \([0,1]\) we therefore have

\begin{equation*} |f|^p \le \|f\|_{\infty}^{\,p} \, \chi_{[0,1]} + |f|^p \chi_N , \end{equation*}

so by order preservation (3.8), additivity of integration (3.16), and 3.15,

\begin{equation*} \int |f|^p \, d\lambda \le \|f\|_{\infty}^{\,p} \, \lambda([0,1]) + \int |f|^p \chi_N \, d\lambda = \|f\|_{\infty}^{\,p} + \int |f|^p \chi_N \, d\lambda . \end{equation*}

The last integral equals \(0\). Indeed, let \(g = |f|^p \chi_N\), a nonnegative measurable function vanishing off \(N\), and let \(A_1, \ldots, A_m\) be any partition of \([0,1]\) into Borel sets. Each term \(\lambda(A_j) \inf_{A_j} g\) of the corresponding lower Lebesgue sum (3.2) is \(0\): if \(A_j \subseteq N\) then \(\lambda(A_j) = 0\) (and \(0 \cdot \infty = 0\) by the book’s convention), while if \(A_j \not\subseteq N\) then \(A_j\) contains a point where \(g\) vanishes, so \(\inf_{A_j} g = 0\). Hence every lower Lebesgue sum is \(0\), and the integral, being the supremum of these sums (3.3), is \(0\) as well. Therefore \(\int |f|^p \, d\lambda \le \|f\|_{\infty}^{\,p} < \infty\), so \(f \in \mathcal{L}^p([0,1])\).

So it suffices to produce a function that lies in \(\mathcal{L}^p([0,1])\) for every \(p < \infty\) but not in \(\mathcal{L}^{\infty}([0,1])\). For \(n \in \mathbb{Z}^{+}\) let

\begin{equation*} A_n = \bigl( 2^{-n}, \, 2^{-n+1} \bigr] , \end{equation*}

so the sets \(A_1, A_2, \ldots\) are disjoint Borel subsets of \([0,1]\) with \(\bigcup_{n=1}^{\infty} A_n = (0,1]\) and \(\lambda(A_n) = 2^{-n+1} - 2^{-n} = 2^{-n}\). Define \(f : [0,1] \to \mathbb{R}\) by

\begin{equation*} f(x) = \sum_{n=1}^{\infty} n \, \chi_{A_n}(x) \end{equation*}

(so \(f(x) = n\) for \(x \in A_n\), and \(f(0) = 0\)). The function \(f\) is Borel measurable because for each \(t \in \mathbb{R}\) the set \(f^{-1}\bigl((t, \infty)\bigr)\) is the union of those \(A_n\) with \(n > t\), together with \(\{0\}\) in case \(t < 0\); this is a Borel set, and the intervals \((t, \infty)\) generate the Borel subsets of \(\mathbb{R}\).

The function \(f\) is not in \(\mathcal{L}^{\infty}([0,1])\): if \(t > 0\) and \(n\) is an integer with \(n > t\), then

\begin{equation*} \lambda\bigl( \{ x : |f(x)| > t \} \bigr) \ge \lambda(A_n) = 2^{-n} > 0 . \end{equation*}

Thus no \(t > 0\) belongs to the set whose infimum defines \(\|f\|_{\infty}\) in 7.1; that set is empty, so \(\|f\|_{\infty} = \infty\) and \(f \notin \mathcal{L}^{\infty}([0,1])\).

The function \(f\) is in \(\mathcal{L}^p([0,1])\) for every \(p \in (0, \infty)\). Fix such a \(p\). Because the sets \(A_n\) are disjoint,

\begin{equation*} |f|^p = \sum_{n=1}^{\infty} n^p \chi_{A_n} , \end{equation*}

and the partial sums \(g_N = \sum_{n=1}^{N} n^p \chi_{A_n}\) form an increasing sequence of nonnegative simple functions converging pointwise to \(|f|^p\). By 3.15,

\begin{equation*} \int g_N \, d\lambda = \sum_{n=1}^{N} n^p \lambda(A_n) = \sum_{n=1}^{N} \frac{n^p}{2^n} , \end{equation*}

so the Monotone Convergence Theorem (3.11) gives

\begin{equation*} \int |f|^p \, d\lambda = \sum_{n=1}^{\infty} \frac{n^p}{2^n} . \end{equation*}

This series converges: writing \(c_n = n^p 2^{-n}\), we have

\begin{equation*} \frac{c_{n+1}}{c_n} = \frac{1}{2}\Bigl(1 + \frac{1}{n}\Bigr)^{p} \longrightarrow \frac{1}{2} \end{equation*}

as \(n \to \infty\), so there exists \(N_0\) with \(c_{n+1} \le \frac{3}{4} c_n\) for all \(n \ge N_0\); hence \(c_n \le c_{N_0} (3/4)^{n - N_0}\) for \(n \ge N_0\), and comparison with a convergent geometric series shows \(\sum_{n=1}^{\infty} c_n < \infty\). Therefore \(\|f\|_p < \infty\), so \(f \in \mathcal{L}^p([0,1])\).

Thus \(f \in \bigcap_{p<\infty} \mathcal{L}^p([0,1])\) but \(f \notin \mathcal{L}^{\infty}([0,1])\), and consequently

\begin{equation*} \bigcap_{p < \infty} \mathcal{L}^p([0,1]) \ne \mathcal{L}^{\infty}([0,1]) , \end{equation*}

the containment \(\mathcal{L}^\infty([0,1]) \subseteq \bigcap_{p<\infty} \mathcal{L}^p([0,1])\) being strict.

Problem (7A.13)

Show that

\begin{equation*} \bigcup_{p > 1} \mathcal{L}^p([0,1]) \ne \mathcal{L}^1([0,1]) . \end{equation*}

Solution

Again \(\lambda\) is Lebesgue measure on \([0,1]\), a finite measure with \(\lambda([0,1]) = 1\).

First, \(\bigcup_{p>1} \mathcal{L}^p([0,1]) \subseteq \mathcal{L}^1([0,1])\). Indeed, if \(1 < p < \infty\) then 7.10 (applied with the exponents \(1 < p\) and the finite measure \(\lambda_{[0,1]}\)) gives \(\mathcal{L}^p([0,1]) \subseteq \mathcal{L}^1([0,1])\); and if \(p = \infty\) then \(\mathcal{L}^{\infty}([0,1]) \subseteq \mathcal{L}^1([0,1])\), as shown at the start of the solution to Exercise 12.

So it suffices to produce \(f \in \mathcal{L}^1([0,1])\) that lies in no \(\mathcal{L}^p([0,1])\) with \(p > 1\). As in Exercise 12, for \(n \in \mathbb{Z}^{+}\) let

\begin{equation*} A_n = \bigl(2^{-n}, \, 2^{-n+1}\bigr] , \end{equation*}

so that \(A_1, A_2, \ldots\) are disjoint Borel sets with \(\lambda(A_n) = 2^{-n}\), and define the Borel measurable function

\begin{equation*} f = \sum_{n=1}^{\infty} \frac{2^n}{n^2} \, \chi_{A_n} \end{equation*}

(and \(f(0) = 0\)); measurability holds because \(f^{-1}((t,\infty))\) is, for each \(t \in \mathbb{R}\), the union of those \(A_n\) with \(2^n/n^2 > t\), together with \(\{0\}\) in case \(t < 0\).

We record a growth fact used twice below: if \(a > 0\) and \(s > 0\), then

\begin{equation*} \frac{2^{na}}{n^s} \longrightarrow \infty \quad \text{as } n \to \infty . \end{equation*}

To see this, pick an integer \(m > s\). From the power series for the exponential function, \(e^t \ge t^m / m!\) for all \(t \ge 0\); taking \(t = n a \ln 2\) gives \(2^{na} \ge (a \ln 2)^m n^m / m!\), so \(\frac{2^{na}}{n^s} \ge \frac{(a \ln 2)^m}{m!} \, n^{m - s} \to \infty\).

Now compute norms. Exactly as in Exercise 12 (disjointness of the \(A_n\), the formula 3.15 for integrals of linear combinations of characteristic functions, and the Monotone Convergence Theorem 3.11 applied to the increasing partial sums), for \(0 < p < \infty\) we get

\begin{equation*} \int |f|^p \, d\lambda = \sum_{n=1}^{\infty} \Bigl(\frac{2^n}{n^2}\Bigr)^{p} \lambda(A_n) = \sum_{n=1}^{\infty} \frac{2^{n(p-1)}}{n^{2p}} . \end{equation*}

Taking \(p = 1\) gives

\begin{equation*} \|f\|_1 = \int |f| \, d\lambda = \sum_{n=1}^{\infty} \frac{1}{n^2} < \infty , \end{equation*}

the convergence holding because \(\frac{1}{n^2} \le \frac{1}{n(n-1)} = \frac{1}{n-1} - \frac{1}{n}\) for \(n \ge 2\), so the partial sums are bounded by \(2\). Hence \(f \in \mathcal{L}^1([0,1])\).

Now suppose \(1 < p < \infty\). Then \(p - 1 > 0\), so by the growth fact above with \(a = p-1\) and \(s = 2p\),

\begin{equation*} \frac{2^{n(p-1)}}{n^{2p}} \longrightarrow \infty \quad \text{as } n \to \infty . \end{equation*}

The terms of the series do not tend to \(0\), so the series diverges and \(\int |f|^p \, d\lambda = \infty\). Thus \(f \notin \mathcal{L}^p([0,1])\).

Finally, \(f \notin \mathcal{L}^{\infty}([0,1])\): by the growth fact with \(a = 1\) and \(s = 2\), we have \(2^n/n^2 \to \infty\), so for every \(t > 0\) there is an \(n\) with \(2^n/n^2 > t\), whence \(\lambda(\{x : |f(x)| > t\}) \ge \lambda(A_n) > 0\); as in Exercise 12 this gives \(\|f\|_{\infty} = \infty\).

Therefore \(f \in \mathcal{L}^1([0,1])\) but \(f \notin \mathcal{L}^p([0,1])\) for any \(p > 1\), so

\begin{equation*} f \in \mathcal{L}^1([0,1]) \setminus \bigcup_{p>1} \mathcal{L}^p([0,1]) , \end{equation*}

and hence \(\bigcup_{p>1} \mathcal{L}^p([0,1]) \ne \mathcal{L}^1([0,1])\), the containment proved at the outset being strict.

Problem (7A.14)

Suppose \(p, q \in (0, \infty]\), with \(p \ne q\). Prove that neither of the sets \(\mathcal{L}^p(\mathbb{R})\) and \(\mathcal{L}^q(\mathbb{R})\) is a subset of the other.

Solution

By symmetry we may assume \(p < q\); thus \(p \in (0, \infty)\) and \(q \in (0, \infty]\). Here \(\lambda\) is Lebesgue measure on \(\mathbb{R}\) and \(\mathcal{L}^p(\mathbb{R}) = \mathcal{L}^p(\lambda)\) as in 7.11. We must produce

\begin{equation*} f \in \mathcal{L}^p(\mathbb{R}) \setminus \mathcal{L}^q(\mathbb{R}) \qquad \text{and} \qquad g \in \mathcal{L}^q(\mathbb{R}) \setminus \mathcal{L}^p(\mathbb{R}) . \end{equation*}

The first function will concentrate near \(0\) (a local blow-up hurts the larger exponent), the second will spread out near \(\infty\) (a heavy tail hurts the smaller exponent).

Throughout we use the following computation. If \(E_1, E_2, \ldots\) are disjoint Borel subsets of \(\mathbb{R}\) and \(c_1, c_2, \ldots \in [0, \infty)\), then \(h = \sum_{n=1}^{\infty} c_n \chi_{E_n}\) is Borel measurable (for each \(t \ge 0\), \(h^{-1}((t,\infty))\) is the union of those \(E_n\) with \(c_n > t\), while for \(t < 0\) it is all of \(\mathbb{R}\), since \(h \ge 0\)), and for \(0 < s < \infty\) the increasing partial sums \(\sum_{n=1}^{N} c_n^{\,s} \chi_{E_n}\) converge pointwise to \(|h|^s\), so 3.15 and the Monotone Convergence Theorem (3.11) give

\begin{equation*} \int |h|^s \, d\lambda = \sum_{n=1}^{\infty} c_n^{\,s} \, \lambda(E_n) . \end{equation*}

Moreover, if \(\sup_n c_n = \infty\) and every \(\lambda(E_n) > 0\), then \(\|h\|_{\infty} = \infty\), since for each \(t > 0\) there is an \(n\) with \(c_n > t\) and then \(\lambda(\{x : |h(x)| > t\}) \ge \lambda(E_n) > 0\); while if \(\sup_n c_n < \infty\) then clearly \(\|h\|_{\infty} \le \sup_n c_n < \infty\).

A function in \(\mathcal{L}^p(\mathbb{R})\) but not in \(\mathcal{L}^q(\mathbb{R})\). Choose a real number \(r\) with \(p < r \le q\): take \(r = q\) if \(q < \infty\), and \(r = p + 1\) if \(q = \infty\). For \(n \in \mathbb{Z}^{+}\) let

\begin{equation*} A_n = \bigl(2^{-n},\, 2^{-n+1}\bigr] , \qquad \lambda(A_n) = 2^{-n} , \end{equation*}

these sets being disjoint with union \((0,1]\), and define

\begin{equation*} f = \sum_{n=1}^{\infty} 2^{n/r} \chi_{A_n} \end{equation*}

(so \(f = 0\) outside \((0,1]\)). For \(0 < s < \infty\) the computation above gives

\begin{equation*} \int |f|^s \, d\lambda = \sum_{n=1}^{\infty} 2^{ns/r} \, 2^{-n} = \sum_{n=1}^{\infty} \bigl(2^{\,s/r - 1}\bigr)^{n} , \end{equation*}

a geometric series with ratio \(2^{s/r - 1}\). Taking \(s = p\): since \(p < r\) we have \(2^{p/r - 1} < 1\), so the series converges and \(\|f\|_p < \infty\); thus \(f \in \mathcal{L}^p(\mathbb{R})\). If \(q < \infty\), then \(r = q\) and taking \(s = q\) gives ratio \(2^0 = 1\), so

\begin{equation*} \int |f|^q \, d\lambda = \sum_{n=1}^{\infty} 1 = \infty , \end{equation*}

and \(f \notin \mathcal{L}^q(\mathbb{R})\). If \(q = \infty\), then \(\sup_n 2^{n/r} = \infty\) and each \(\lambda(A_n) > 0\), so \(\|f\|_{\infty} = \infty\) and again \(f \notin \mathcal{L}^q(\mathbb{R})\).

A function in \(\mathcal{L}^q(\mathbb{R})\) but not in \(\mathcal{L}^p(\mathbb{R})\). For \(n \in \mathbb{Z}^{+}\) let

\begin{equation*} B_n = (n, \, n+1] , \qquad \lambda(B_n) = 1 , \end{equation*}

these sets being disjoint, and define

\begin{equation*} g = \sum_{n=1}^{\infty} n^{-1/p} \chi_{B_n} \end{equation*}

(so \(g = 0\) outside \((1, \infty)\)); this makes sense because \(p < \infty\). For \(0 < s < \infty\) the computation above gives

\begin{equation*} \int |g|^s \, d\lambda = \sum_{n=1}^{\infty} n^{-s/p} . \end{equation*}

Taking \(s = p\) gives \(\int |g|^p \, d\lambda = \sum_{n=1}^{\infty} \frac{1}{n} = \infty\), the harmonic series being divergent (as shown in the solution to Exercise 11). Hence \(g \notin \mathcal{L}^p(\mathbb{R})\).

If \(q < \infty\), take \(s = q\): then \(\frac{q}{p} > 1\), and \(\sum_{n=1}^{\infty} n^{-q/p} < \infty\) by the estimate for \(\sum n^{-t}\) with \(t > 1\) proved in the solution to Exercise 11. Hence \(\|g\|_q < \infty\) and \(g \in \mathcal{L}^q(\mathbb{R})\). If \(q = \infty\), then \(\sup_n n^{-1/p} = 1 < \infty\), so \(\|g\|_{\infty} \le 1\) and again \(g \in \mathcal{L}^q(\mathbb{R})\).

Conclusion. We have produced \(f \in \mathcal{L}^p(\mathbb{R})\) with \(f \notin \mathcal{L}^q(\mathbb{R})\), which shows \(\mathcal{L}^p(\mathbb{R}) \not\subseteq \mathcal{L}^q(\mathbb{R})\), and \(g \in \mathcal{L}^q(\mathbb{R})\) with \(g \notin \mathcal{L}^p(\mathbb{R})\), which shows \(\mathcal{L}^q(\mathbb{R}) \not\subseteq \mathcal{L}^p(\mathbb{R})\). Thus neither of \(\mathcal{L}^p(\mathbb{R})\) and \(\mathcal{L}^q(\mathbb{R})\) is a subset of the other.

Note the contrast with 7.10 and with Exercise 10: on a finite measure space such as \([0,1]\) the larger exponent gives the smaller space, and for counting measure on \(\mathbb{Z}^{+}\) the smaller exponent gives the smaller space. Lebesgue measure on \(\mathbb{R}\) has both sets of small measure and sets of large measure, so neither containment survives.

Problem (7A.15)

Show that there exists \(f \in \mathcal{L}^2(\mathbb{R})\) such that \(f \notin \mathcal{L}^p(\mathbb{R})\) for all \(p \in (0, \infty] \setminus \{2\}\).

Solution

Define \(f \colon \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} f(x) = \begin{cases} \dfrac{1}{\sqrt{x}\,\lvert \ln x \rvert} & \text{if } x \in \left(0, \tfrac{1}{2}\right) \cup (2, \infty),\\[2mm] 0 & \text{otherwise.}\end{cases} \end{equation*}

Then \(f\) is continuous on the open set \(\left(0,\tfrac12\right) \cup (2,\infty)\) and vanishes off that set, so \(f\) is a Borel measurable function.

Note that \(\ln x \neq 0\) on the two intervals above, so \(f\) is well defined, and that \(\lvert \ln x\rvert = \ln x\) when \(x > 2\).

Step 1: \(f \in \mathcal{L}^2(\mathbb{R})\). The substitution \(u = \ln x\), \(du = dx/x\), gives

\begin{equation*} \int_0^{1/2} \frac{dx}{x (\ln x)^2} = \int_{-\infty}^{-\ln 2} \frac{du}{u^2} = \frac{1}{\ln 2} \qquad\text{and}\qquad \int_2^{\infty} \frac{dx}{x (\ln x)^2} = \int_{\ln 2}^{\infty} \frac{du}{u^2} = \frac{1}{\ln 2}. \end{equation*}

Because \(f\) is continuous and nonnegative on each of the two intervals, the Lebesgue integral of \(f^2\) over each interval agrees with the improper Riemann integral computed above. Hence

\begin{equation*} \|f\|_2^{\,2} = \int \lvert f\rvert^2 \, d\lambda = \frac{2}{\ln 2} < \infty , \end{equation*}

so \(f \in \mathcal{L}^2(\mathbb{R})\).

Step 2: \(f \notin \mathcal{L}^p(\mathbb{R})\) for \(0 < p < 2\). Fix such a \(p\) and put \(q = p/2 \in (0,1)\). On \((2,\infty)\) we have

\begin{equation*} \lvert f(x)\rvert^p = \frac{1}{x^{q}\,(\ln x)^{p}} . \end{equation*}

Because \((\ln x)^p \big/ x^{(1-q)/2} \to 0\) as \(x \to \infty\), there exists \(M \geq 2\) such that \((\ln x)^p \leq x^{(1-q)/2}\) for all \(x \geq M\). Thus for \(x \geq M\),

\begin{equation*} \lvert f(x)\rvert^p \geq \frac{1}{x^{q} \cdot x^{(1-q)/2}} = \frac{1}{x^{(1+q)/2}} . \end{equation*}

Since \(q < 1\), we have \((1+q)/2 < 1\), and therefore \(\int_M^{\infty} x^{-(1+q)/2}\,dx = \infty\). Hence \(\int \lvert f\rvert^p \, d\lambda = \infty\) and \(f \notin \mathcal{L}^p(\mathbb{R})\).

Step 3: \(f \notin \mathcal{L}^p(\mathbb{R})\) for \(2 < p < \infty\). Fix such a \(p\) and put \(q = p/2 > 1\); choose \(r\) with \(1 < r < q\). On \(\left(0,\tfrac12\right)\) we have

\begin{equation*} \lvert f(x)\rvert^p = \frac{1}{x^{q}\,\lvert \ln x\rvert^{p}} . \end{equation*}

Because \(x^{q - r} \lvert \ln x \rvert^{p} \to 0\) as \(x \downarrow 0\) (a positive power of \(x\) beats any power of \(\lvert \ln x\rvert\)), there exists \(\delta \in \left(0, \tfrac12\right)\) such that \(\lvert \ln x\rvert^{p} \leq x^{-(q-r)}\) for all \(x \in (0,\delta)\). Thus for such \(x\),

\begin{equation*} \lvert f(x)\rvert^p \geq \frac{x^{q-r}}{x^{q}} = \frac{1}{x^{r}} . \end{equation*}

Since \(r > 1\), we have \(\int_0^{\delta} x^{-r}\,dx = \infty\). Hence \(\int \lvert f\rvert^p \, d\lambda = \infty\) and \(f \notin \mathcal{L}^p(\mathbb{R})\).

Step 4: \(f \notin \mathcal{L}^{\infty}(\mathbb{R})\). As \(x \downarrow 0\) we have \(\sqrt{x}\,\lvert \ln x\rvert \to 0\), so \(f(x) \to \infty\). Hence for each \(t > 0\) there exists \(\delta_t \in \left(0,\tfrac12\right)\) such that \(f(x) > t\) for all \(x \in (0, \delta_t)\); thus

\begin{equation*} \lambda\bigl(\{x \in \mathbb{R} : \lvert f(x)\rvert > t\}\bigr) \geq \delta_t > 0 \end{equation*}

for every \(t > 0\). Therefore the set appearing in the definition 7.1 of the essential supremum is empty, so \(\|f\|_{\infty} = \inf \emptyset = \infty\) and \(f \notin \mathcal{L}^{\infty}(\mathbb{R})\).

Combining Steps 1 through 4: \(f \in \mathcal{L}^2(\mathbb{R})\) and \(f \notin \mathcal{L}^p(\mathbb{R})\) for every \(p \in (0,\infty] \setminus \{2\}\).

Problem (7A.16)

Suppose \((X, \mathcal{S}, \mu)\) is a finite measure space. Prove that

\begin{equation*} \lim_{p \to \infty} \|f\|_p = \|f\|_{\infty} \end{equation*}

for every \(\mathcal{S}\)-measurable function \(f \colon X \to \mathbf{F}\).

Solution

Throughout, \(\|f\|_p\) and \(\|f\|_{\infty}\) are as in 7.1; no integrability of \(f\) is assumed, and the values \(\|f\|_p = \infty\) are allowed.

The degenerate case. If \(\mu(X) = 0\), then every set in \(\mathcal{S}\) has measure \(0\), so \(\mu(\{x \in X : \lvert f(x)\rvert > t\}) = 0\) for every \(t > 0\); thus every positive number belongs to the set whose infimum defines \(\|f\|_{\infty}\) in 7.1, and hence \(\|f\|_{\infty} = 0\). Also \(\int \lvert f\rvert^p \, d\mu = 0\) for every \(p \in (0,\infty)\), so \(\|f\|_p = 0\). Thus the limit holds. Assume from now on that \(\mu(X) > 0\).

Step 1: \(\lvert f \rvert \leq \|f\|_{\infty}\) almost everywhere. Suppose \(\|f\|_{\infty} < \infty\). For each \(n \in \mathbb{Z}^+\), the definition of \(\|f\|_{\infty}\) as an infimum produces \(t_n\) with \(\|f\|_{\infty} \leq t_n < \|f\|_{\infty} + \tfrac1n\) and \(\mu(\{\lvert f\rvert > t_n\}) = 0\). Because

\begin{equation*} \{x \in X : \lvert f(x)\rvert > \|f\|_{\infty}\} = \bigcup_{n=1}^{\infty} \Bigl\{x \in X : \lvert f(x)\rvert > \|f\|_{\infty} + \tfrac1n \Bigr\} \subseteq \bigcup_{n=1}^{\infty}\{x \in X : \lvert f(x)\rvert > t_n\}, \end{equation*}

countable subadditivity of \(\mu\) gives \(\mu(\{\lvert f\rvert > \|f\|_{\infty}\}) = 0\), as claimed.

Step 2: \(\limsup_{p \to \infty} \|f\|_p \leq \|f\|_{\infty}\). If \(\|f\|_{\infty} = \infty\) this is trivial, so suppose \(\|f\|_{\infty} < \infty\). By Step 1, \(\lvert f\rvert^p \leq \|f\|_{\infty}^{\,p}\) almost everywhere, and hence for every \(p \in (0,\infty)\)

\begin{equation*} \int \lvert f\rvert^p \, d\mu \leq \|f\|_{\infty}^{\,p}\, \mu(X), \end{equation*}

using \(\mu(X) < \infty\). Raising to the power \(1/p\) gives

\begin{equation*} \|f\|_p \leq \|f\|_{\infty} \,\mu(X)^{1/p} . \end{equation*}

Because \(0 < \mu(X) < \infty\), we have \(\mu(X)^{1/p} = e^{(\ln \mu(X))/p} \to 1\) as \(p \to \infty\). Therefore \(\limsup_{p \to \infty} \|f\|_p \leq \|f\|_{\infty}\).

Step 3: \(\liminf_{p \to \infty} \|f\|_p \geq \|f\|_{\infty}\). If \(\|f\|_{\infty} = 0\) this is trivial because \(\|f\|_p \geq 0\). So suppose \(\|f\|_{\infty} > 0\) and let \(t\) be any number with \(0 < t < \|f\|_{\infty}\). Since \(t\) is smaller than the infimum defining \(\|f\|_{\infty}\), the number \(t\) does not belong to the set \(\{s > 0 : \mu(\{\lvert f\rvert > s\}) = 0\}\); thus

\begin{equation*} c := \mu\bigl(\{x \in X : \lvert f(x)\rvert > t\}\bigr) > 0 , \end{equation*}

and \(c \leq \mu(X) < \infty\). Writing \(E = \{x \in X : \lvert f(x)\rvert > t\}\), for every \(p \in (0,\infty)\) we get

\begin{equation*} \int \lvert f \rvert^p \, d\mu \geq \int_E \lvert f\rvert^p \, d\mu \geq t^p \mu(E) = t^p c , \end{equation*}

hence \(\|f\|_p \geq t\, c^{1/p}\). Because \(0 < c < \infty\), we have \(c^{1/p} \to 1\) as \(p \to \infty\), and therefore

\begin{equation*} \liminf_{p \to \infty} \|f\|_p \geq t . \end{equation*}

Letting \(t\) increase to \(\|f\|_{\infty}\) (through numbers less than \(\|f\|_{\infty}\), which is legitimate whether \(\|f\|_{\infty}\) is finite or infinite) gives \(\liminf_{p \to \infty} \|f\|_p \geq \|f\|_{\infty}\).

Steps 2 and 3 together show that \(\lim_{p \to \infty}\|f\|_p\) exists in \([0,\infty]\) and equals \(\|f\|_{\infty}\).

Problem (7A.17)

Suppose \(\mu\) is a measure, \(0 < p \leq \infty\), and \(f \in \mathcal{L}^p(\mu)\). Prove that for every \(\varepsilon > 0\), there exists a simple function \(g \in \mathcal{L}^p(\mu)\) such that \(\|f - g\|_p < \varepsilon\).

[This exercise extends 3.44.]

Solution

Let \((X, \mathcal{S}, \mu)\) be the underlying measure space. Recall from 2.88 that a simple function is one taking on only finitely many values; the functions we produce will all be \(\mathcal{S}\)-measurable.

A pointwise approximation. Write \(u = \operatorname{Re} f\) and \(v = \operatorname{Im} f\) (with \(v = 0\) if \(\mathbf{F} = \mathbb{R}\)); these are \(\mathcal{S}\)-measurable real-valued functions. By 2.89 there exist simple \(\mathcal{S}\)-measurable functions \(u_1, u_2, \ldots\) and \(v_1, v_2, \ldots\) from \(X\) to \(\mathbb{R}\) such that

\begin{equation*} \lvert u_k \rvert \leq \lvert u\rvert, \quad \lvert v_k\rvert \leq \lvert v \rvert, \quad u_k \to u \ \text{ and } \ v_k \to v \ \text{ pointwise on } X, \end{equation*}

and such that the convergence is uniform on \(X\) when the function being approximated is bounded [2.89(b), (c), (d)]. Put \(g_k = u_k + i v_k\). Each \(g_k\) takes on only finitely many values and is \(\mathcal{S}\)-measurable, so each \(g_k\) is a simple \(\mathcal{S}\)-measurable function; moreover

\begin{equation*} \lvert g_k \rvert \leq \lvert u_k\rvert + \lvert v_k\rvert \leq \lvert u \rvert + \lvert v\rvert \leq 2 \lvert f \rvert \qquad\text{and}\qquad g_k(x) \to f(x) \ \text{ for every } x \in X . \end{equation*}

In particular \(\|g_k\|_p \leq 2\|f\|_p < \infty\), so \(g_k \in \mathcal{L}^p(\mu)\) for every \(p \in (0,\infty]\).

Case 1: \(0 < p < \infty\). Set \(h_k = \lvert f - g_k\rvert^p\). Then each \(h_k\) is \(\mathcal{S}\)-measurable, \(h_k(x) \to 0\) for every \(x \in X\), and

\begin{equation*} h_k = \lvert f - g_k \rvert^p \leq \bigl(\lvert f \rvert + \lvert g_k\rvert\bigr)^p \leq \bigl(3\lvert f\rvert\bigr)^p = 3^p \lvert f\rvert^p . \end{equation*}

Because \(f \in \mathcal{L}^p(\mu)\), the dominating function \(3^p\lvert f\rvert^p\) has finite integral. Hence the Dominated Convergence Theorem (3.31) gives

\begin{equation*} \lim_{k \to \infty} \int \lvert f - g_k\rvert^p \, d\mu = 0 , \end{equation*}

that is, \(\|f - g_k\|_p \to 0\). Given \(\varepsilon > 0\), choose \(k\) with \(\|f - g_k\|_p < \varepsilon\) and take \(g = g_k\), a simple function in \(\mathcal{L}^p(\mu)\).

Case 2: \(p = \infty\). Let \(M = \|f\|_{\infty} < \infty\) and let \(E = \{x \in X : \lvert f(x)\rvert \leq M\}\). Exactly as in Step 1 of the solution to 7A.16, \(\mu(X \setminus E) = 0\). Define \(\tilde f = f \chi_E\); then \(\tilde f\) is \(\mathcal{S}\)-measurable, \(\lvert \tilde f \rvert \leq M\) everywhere on \(X\), and \(\tilde f = f\) outside a set of measure \(0\), so \(\|f - \tilde f\|_{\infty} = 0\).

Now \(\operatorname{Re} \tilde f\) and \(\operatorname{Im} \tilde f\) are bounded \(\mathcal{S}\)-measurable real-valued functions, so 2.89(d) provides simple \(\mathcal{S}\)-measurable functions \(a, b \colon X \to \mathbb{R}\) with

\begin{equation*} \sup_{x \in X} \lvert \operatorname{Re}\tilde f(x) - a(x)\rvert < \frac{\varepsilon}{4} \qquad\text{and}\qquad \sup_{x \in X}\lvert \operatorname{Im}\tilde f(x) - b(x)\rvert < \frac{\varepsilon}{4} . \end{equation*}

Let \(g = a + ib\) (take \(g = a\) if \(\mathbf{F} = \mathbb{R}\)). Then \(g\) takes on only finitely many values and is \(\mathcal{S}\)-measurable, so \(g\) is a simple function, and \(\lvert \tilde f(x) - g(x)\rvert < \varepsilon/2\) for every \(x \in X\). Consequently \(\lvert g \rvert \leq M + \varepsilon/2\) everywhere, so \(\|g\|_{\infty} < \infty\) and \(g \in \mathcal{L}^{\infty}(\mu)\). Finally, using the triangle inequality for \(\|\cdot\|_{\infty}\) (Exercise 1 in this section),

\begin{equation*} \|f - g\|_{\infty} \leq \|f - \tilde f\|_{\infty} + \|\tilde f - g\|_{\infty} \leq 0 + \frac{\varepsilon}{2} < \varepsilon . \end{equation*}

In both cases we have produced a simple function \(g \in \mathcal{L}^p(\mu)\) with \(\|f - g\|_p < \varepsilon\), as desired.

Problem (7A.18)

Suppose \(0 < p < \infty\) and \(f \in \mathcal{L}^p(\mathbb{R})\). Prove that for every \(\varepsilon > 0\), there exists a step function \(g \in \mathcal{L}^p(\mathbb{R})\) such that \(\|f - g\|_p < \varepsilon\).

[This exercise extends 3.47.]

Solution

Here \(\lambda\) denotes Lebesgue measure and \(\lvert A \rvert = \lambda(A)\). As in 3.46, a step function is a function of the form \(a_1 \chi_{I_1} + \cdots + a_n \chi_{I_n}\), where \(I_1, \ldots, I_n\) are intervals of \(\mathbb{R}\) and \(a_1, \ldots, a_n\) are nonzero; when \(\mathbf{F} = \mathbb{C}\) we allow the coefficients \(a_k\) to be complex.

A substitute for the triangle inequality when \(p < 1\). Put \(q = \min\{p, 1\}\). We claim that

\begin{equation*} \Bigl\| \sum_{k=1}^m f_k \Bigr\|_p^{\,q} \leq \sum_{k=1}^m \|f_k\|_p^{\,q} \end{equation*}

for all \(f_1, \ldots, f_m \in \mathcal{L}^p(\mu)\) and every measure \(\mu\). If \(p \geq 1\), then \(q = 1\) and this is Minkowski’s inequality (7.14) together with induction on \(m\). If \(0 < p < 1\), then \(q = p\), and we first note that

\begin{equation*} (a + b)^p \leq a^p + b^p \qquad \text{for all } a, b \geq 0 : \end{equation*}

this is clear if \(a + b = 0\), and otherwise, writing \(s = a/(a+b)\) and \(t = b/(a+b)\), we have \(s, t \in [0,1]\) and \(s + t = 1\), so \(s^p \geq s\) and \(t^p \geq t\) (because \(r^p \geq r\) for \(r \in [0,1]\) and \(p \leq 1\)), giving \(a^p + b^p = (a+b)^p (s^p + t^p) \geq (a+b)^p\). Iterating this inequality and integrating,

\begin{equation*} \Bigl\| \sum_{k=1}^m f_k \Bigr\|_p^{\,p} = \int \Bigl\lvert \sum_{k=1}^m f_k \Bigr\rvert^p d\mu \leq \int \Bigl( \sum_{k=1}^m \lvert f_k \rvert \Bigr)^{p} d\mu \leq \sum_{k=1}^m \int \lvert f_k\rvert^p \, d\mu = \sum_{k=1}^m \|f_k\|_p^{\,p} , \end{equation*}

which proves the claim.

Approximating a measurable set of finite measure. Suppose \(A \subseteq \mathbb{R}\) is Lebesgue measurable with \(\lvert A\rvert < \infty\) and \(\delta > 0\). By part (e) of 2.71 there is an open set \(G \supseteq A\) with \(\lvert G \setminus A\rvert < \delta/2\); thus \(\lvert G\rvert \leq \lvert A\rvert + \delta/2 < \infty\). Every open subset of \(\mathbb{R}\) is a countable union of disjoint open intervals (this fact is used in the proof of 3.47), so write \(G = \bigcup_{j=1}^{\infty} J_j\) with \(J_1, J_2, \ldots\) disjoint open intervals; each \(J_j\) is bounded because \(\lvert G \rvert < \infty\). Since \(\sum_{j=1}^{\infty} \lvert J_j\rvert = \lvert G\rvert < \infty\), there exists \(N\) with \(\sum_{j > N} \lvert J_j\rvert < \delta/2\). Let \(E = J_1 \cup \cdots \cup J_N\), a union of finitely many disjoint bounded open intervals. Then

\begin{equation*} A \setminus E \subseteq G \setminus E = \bigcup_{j > N} J_j \quad\text{and}\quad E \setminus A \subseteq G \setminus A , \end{equation*}

so \(\lvert A \setminus E\rvert + \lvert E \setminus A \rvert < \delta\). Since \(\lvert \chi_A - \chi_E\rvert = \chi_{(A \setminus E) \cup (E \setminus A)}\), this gives

\begin{equation*} \|\chi_A - \chi_E\|_p^{\,p} = \lvert A \setminus E\rvert + \lvert E \setminus A\rvert < \delta . \end{equation*}

Proof of the exercise. Let \(\varepsilon > 0\). By 7A.17 there is a simple function \(h \in \mathcal{L}^p(\mathbb{R})\) with

\begin{equation*} \|f - h\|_p^{\,q} < \frac{\varepsilon^q}{2} . \end{equation*}

If \(h = 0\), then \(g = 0\) is a step function (see the remark following 3.46) and we are done. Otherwise write \(h\) in its standard representation

\begin{equation*} h = a_1 \chi_{A_1} + \cdots + a_n \chi_{A_n}, \end{equation*}

where \(a_1, \ldots, a_n\) are the distinct nonzero values of \(h\) and \(A_1, \ldots, A_n\) are disjoint Lebesgue measurable sets (see the discussion following 2.88). Because the \(A_k\) are disjoint,

\begin{equation*} \|h\|_p^{\,p} = \sum_{k=1}^n \lvert a_k\rvert^p \lvert A_k\rvert < \infty , \end{equation*}

which forces \(\lvert A_k \rvert < \infty\) for each \(k\).

For each \(k \in \{1, \ldots, n\}\) choose \(\delta_k > 0\) small enough that

\begin{equation*} \lvert a_k \rvert^{q} \, \delta_k^{\,q/p} < \frac{\varepsilon^q}{2n} , \end{equation*}

and then, by the previous paragraph, choose a set \(E_k\) that is a union of finitely many disjoint bounded open intervals with \(\|\chi_{A_k} - \chi_{E_k}\|_p^{\,p} < \delta_k\). Define

\begin{equation*} g = a_1 \chi_{E_1} + \cdots + a_n \chi_{E_n} . \end{equation*}

Expanding each \(\chi_{E_k}\) as the sum of the characteristic functions of the finitely many disjoint bounded open intervals whose union is \(E_k\), we see that \(g\) is a finite linear combination, with nonzero coefficients, of characteristic functions of bounded intervals; thus \(g\) is a step function. Moreover \(g\) is bounded and vanishes outside a bounded set, so \(g \in \mathcal{L}^p(\mathbb{R})\).

By the \(q\)-subadditivity established above and 7.5(b),

\begin{equation*} \|h - g\|_p^{\,q} \leq \sum_{k=1}^n \bigl\| a_k(\chi_{A_k} - \chi_{E_k})\bigr\|_p^{\,q} = \sum_{k=1}^n \lvert a_k\rvert^{q}\, \|\chi_{A_k} - \chi_{E_k}\|_p^{\,q} < \sum_{k=1}^n \lvert a_k \rvert^q \delta_k^{\,q/p} < \frac{\varepsilon^q}{2} . \end{equation*}

Therefore

\begin{equation*} \|f - g\|_p^{\,q} \leq \|f - h\|_p^{\,q} + \|h - g\|_p^{\,q} < \frac{\varepsilon^q}{2} + \frac{\varepsilon^q}{2} = \varepsilon^q , \end{equation*}

which gives \(\|f - g\|_p < \varepsilon\), as desired.

Problem (7A.19)

Suppose \(0 < p < \infty\) and \(f \in \mathcal{L}^p(\mathbb{R})\). Prove that for every \(\varepsilon > 0\), there exists a continuous function \(g \colon \mathbb{R} \to \mathbf{F}\) such that \(\|f - g\|_p < \varepsilon\) and the set \(\{x \in \mathbb{R} : g(x) \neq 0\}\) is bounded.

[This exercise extends 3.48.]

Solution

As in the solution to 7A.18, put \(q = \min\{p, 1\}\) and recall the inequality proved there:

\begin{equation*} \Bigl\| \sum_{k=1}^m f_k\Bigr\|_p^{\,q} \leq \sum_{k=1}^m \|f_k\|_p^{\,q} \qquad \text{for } f_1, \ldots, f_m \in \mathcal{L}^p(\mathbb{R}). \end{equation*}

Let \(\varepsilon > 0\). By 7A.18 there is a step function \(h \in \mathcal{L}^p(\mathbb{R})\) with

\begin{equation*} \|f - h\|_p^{\,q} < \frac{\varepsilon^q}{2} . \end{equation*}

If \(h = 0\), take \(g = 0\) and we are done. Otherwise write

\begin{equation*} h = a_1 \chi_{I_1} + \cdots + a_n \chi_{I_n}, \end{equation*}

where \(a_1, \ldots, a_n\) are nonzero elements of \(\mathbf{F}\) and \(I_1, \ldots, I_n\) are intervals of \(\mathbb{R}\).

We may assume that \(I_1, \ldots, I_n\) are disjoint and bounded. To see this, note that the finitely many endpoints of the original intervals cut \(\mathbb{R}\) into finitely many disjoint pieces: the singletons consisting of the endpoints themselves, the bounded open intervals between consecutive endpoints, and two unbounded open intervals beyond the extreme endpoints. The function \(h\) is constant on each such piece. Discarding the pieces on which \(h\) equals \(0\) and relabeling, we obtain a representation \(h = a_1\chi_{I_1} + \cdots + a_n \chi_{I_n}\) in which \(a_1, \ldots, a_n\) are nonzero and \(I_1, \ldots, I_n\) are disjoint intervals (if no pieces remain, then \(h = 0\), a case already handled). Because the intervals are now disjoint,

\begin{equation*} \|h\|_p^{\,p} = \sum_{k=1}^n \lvert a_k\rvert^p\, \lvert I_k\rvert < \infty , \end{equation*}

which forces \(\lvert I_k \rvert < \infty\), and hence forces \(I_k\) to be bounded, for each \(k\). For each \(k\) let \(b_k \leq c_k\) denote the endpoints of \(I_k\), so that \(\lvert I_k\rvert = c_k - b_k\).

Continuous approximation of \(\chi_{I}\). Let \(I\) be a bounded interval with endpoints \(b \leq c\), and let \(\delta > 0\). Define \(\varphi \colon \mathbb{R} \to \mathbb{R}\) as follows. If \(c - b \leq 2\delta\), let \(\varphi = 0\); then

\begin{equation*} \|\chi_I - \varphi\|_p^{\,p} = \lvert I\rvert = c - b \leq 2\delta . \end{equation*}

If \(c - b > 2\delta\), let

\begin{equation*} \varphi(x) = \begin{cases} 0 & \text{if } x \leq b \text{ or } x \geq c,\\[1mm] \dfrac{x - b}{\delta} & \text{if } b \leq x \leq b + \delta,\\[2mm] 1 & \text{if } b + \delta \leq x \leq c - \delta,\\[1mm] \dfrac{c - x}{\delta} & \text{if } c - \delta \leq x \leq c. \end{cases} \end{equation*}

Because \(c - b > 2\delta\) we have \(b < b + \delta < c - \delta < c\), and the four formulas agree at the junction points \(b\), \(b + \delta\), \(c - \delta\), \(c\); thus \(\varphi\) is continuous. It takes values in \([0,1]\) and vanishes outside \([b,c]\). If \(x < b\) or \(x > c\), then \(\varphi(x) = 0 = \chi_I(x)\); if \(b + \delta < x < c - \delta\), then \(\varphi(x) = 1 = \chi_I(x)\), because \((b+\delta, c-\delta) \subseteq (b,c) \subseteq I\). Hence \(\chi_I - \varphi\) vanishes outside the set \([b, b+\delta] \cup [c - \delta, c]\), which has measure \(2\delta\), and \(\lvert \chi_I - \varphi\rvert \leq 1\) everywhere. Therefore again

\begin{equation*} \|\chi_I - \varphi\|_p^{\,p} \leq 2\delta . \end{equation*}

In both cases \(\varphi\) is continuous, vanishes outside the bounded set \([b, c]\), and satisfies \(\|\chi_I - \varphi\|_p^{\,p} \leq 2\delta\).

Conclusion. Choose \(\delta > 0\) small enough that

\begin{equation*} \lvert a_k\rvert^{q} (2\delta)^{q/p} < \frac{\varepsilon^q}{2n} \qquad \text{for } k = 1, \ldots, n, \end{equation*}

which is possible since \((2\delta)^{q/p} \to 0\) as \(\delta \downarrow 0\). For each \(k\) let \(\varphi_k\) be the continuous function constructed above for the interval \(I_k\) and this \(\delta\), and set

\begin{equation*} g = a_1 \varphi_1 + \cdots + a_n \varphi_n . \end{equation*}

Then \(g \colon \mathbb{R} \to \mathbf{F}\) is continuous, and \(\{x \in \mathbb{R} : g(x) \neq 0\} \subseteq [b_1, c_1] \cup \cdots \cup [b_n, c_n]\), which is a bounded set. Being continuous and vanishing outside a bounded set, \(g\) is bounded and hence \(g \in \mathcal{L}^p(\mathbb{R})\).

By the \(q\)-subadditivity above and 7.5(b),

\begin{equation*} \|h - g\|_p^{\,q} \leq \sum_{k=1}^n \lvert a_k \rvert^{q}\, \|\chi_{I_k} - \varphi_k\|_p^{\,q} \leq \sum_{k=1}^n \lvert a_k\rvert^{q} (2\delta)^{q/p} < \frac{\varepsilon^q}{2} . \end{equation*}

Therefore

\begin{equation*} \|f - g\|_p^{\,q} \leq \|f - h\|_p^{\,q} + \|h - g\|_p^{\,q} < \varepsilon^q , \end{equation*}

so \(\|f - g\|_p < \varepsilon\), completing the proof.

Problem (7A.20)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space, \(1 < p < \infty\), and \(f, g \in \mathcal{L}^p(\mu)\). Prove that Minkowski’s inequality (7.14) is an equality if and only if there exist nonnegative numbers \(a\) and \(b\), not both \(0\), such that

\begin{equation*} a f(x) = b g(x) \end{equation*}

for almost every \(x \in X\).

Solution

We must show that \(\|f + g\|_p = \|f\|_p + \|g\|_p\) if and only if \(af = bg\) almost everywhere for some nonnegative \(a, b\) that are not both \(0\).

Throughout we use two elementary facts about \(\mathbf{F}\) and about \(t \mapsto t^p\).

Fact 1 (equality in the triangle inequality in \(\mathbf{F}\)). For \(z, w \in \mathbf{F}\) we have \(\lvert z + w\rvert = \lvert z\rvert + \lvert w\rvert\) if and only if \(z \overline{w} \in [0,\infty)\). Indeed,

\begin{equation*} \lvert z + w\rvert^2 = \lvert z\rvert^2 + 2 \operatorname{Re}(z \overline{w}) + \lvert w \rvert^2 \quad\text{and}\quad (\lvert z\rvert + \lvert w\rvert)^2 = \lvert z \rvert^2 + 2 \lvert z \overline{w}\rvert + \lvert w\rvert^2 , \end{equation*}

so equality holds if and only if \(\operatorname{Re}(z\overline w) = \lvert z \overline w\rvert\), which happens if and only if \(z \overline w\) is a nonnegative real number. (When \(\mathbf{F} = \mathbb{R}\) this says \(zw \geq 0\).)

Fact 2 (strict convexity of \(t \mapsto t^p\)). Suppose \(p > 1\), \(s \in (0,1)\), and \(\alpha, \beta \in [0,\infty)\) with \(\alpha \neq \beta\). Then

\begin{equation*} \bigl(s\alpha + (1-s)\beta\bigr)^p < s \alpha^p + (1-s)\beta^p . \end{equation*}

If \(\alpha, \beta > 0\) this follows because \(\varphi(t) = t^p\) satisfies \(\varphi^{\prime\prime}(t) = p(p-1)t^{p-2} > 0\) on \((0,\infty)\), so \(\varphi\) is strictly convex there. If \(\alpha = 0 < \beta\), then \((s \cdot 0 + (1-s)\beta)^p = (1-s)^p \beta^p < (1-s)\beta^p\) because \(0 < 1 - s < 1\) and \(p > 1\); the case \(\beta = 0 < \alpha\) is symmetric. The nonstrict inequality \(\bigl(s\alpha + (1-s)\beta\bigr)^p \leq s\alpha^p + (1-s)\beta^p\) holds for all \(\alpha,\beta \geq 0\) (with equality when \(\alpha = \beta\)).

Proof that the condition is sufficient. Suppose \(a, b \geq 0\) are not both \(0\) and \(af = bg\) almost everywhere. If \(b = 0\), then \(a > 0\), so \(f = 0\) almost everywhere; hence \(f + g = g\) almost everywhere and \(\|f\|_p = 0\), so

\begin{equation*} \|f+g\|_p = \|g\|_p = \|f\|_p + \|g\|_p . \end{equation*}

Now suppose \(b > 0\) and put \(c = a/b \geq 0\). Then \(g = cf\) almost everywhere, so \(f + g = (1+c)f\) almost everywhere, and 7.5(b) gives

\begin{equation*} \|f + g\|_p = (1 + c)\|f\|_p = \|f\|_p + c\|f\|_p = \|f\|_p + \|g\|_p . \end{equation*}

In both cases Minkowski’s inequality is an equality.

Proof that the condition is necessary. Suppose \(\|f + g\|_p = \|f\|_p + \|g\|_p\). If \(\|f\|_p = 0\), then \(f = 0\) almost everywhere and \(a = 1\), \(b = 0\) works. If \(\|g\|_p = 0\), then \(g = 0\) almost everywhere and \(a = 0\), \(b = 1\) works. So assume

\begin{equation*} A := \|f\|_p > 0 \qquad\text{and}\qquad B := \|g\|_p > 0 , \end{equation*}

both finite since \(f, g \in \mathcal{L}^p(\mu)\). Put

\begin{equation*} u = \frac{f}{A}, \qquad v = \frac{g}{B}, \qquad s = \frac{A}{A + B} \in (0,1) , \end{equation*}

so that \(\|u\|_p = \|v\|_p = 1\) by 7.5(b) and \(1 - s = B/(A+B)\). Since \(f + g = (A+B)\bigl(s u + (1-s)v\bigr)\), the assumed equality says exactly that

\begin{equation*} \bigl\| s u + (1-s) v \bigr\|_p = 1 . \end{equation*}

Now, pointwise on \(X\),

\begin{equation*} \bigl\lvert su + (1-s)v \bigr\rvert^p \leq \bigl(s\lvert u\rvert + (1-s)\lvert v\rvert\bigr)^p \leq s \lvert u\rvert^p + (1-s)\lvert v\rvert^p , \end{equation*}

the first inequality by the triangle inequality in \(\mathbf{F}\) and monotonicity of \(t \mapsto t^p\), the second by Fact 2. Integrating and using \(\int \lvert u\rvert^p d\mu = \int \lvert v \rvert^p d\mu = 1\),

\begin{equation*} 1 = \int \bigl\lvert su + (1-s)v\bigr\rvert^p \, d\mu \leq \int \bigl(s\lvert u\rvert + (1-s)\lvert v\rvert\bigr)^p d\mu \leq s + (1-s) = 1 . \end{equation*}

Hence both inequalities are equalities. All the functions involved lie in \(\mathcal{L}^1(\mu)\) (each is dominated by \(s\lvert u\rvert^p + (1-s)\lvert v\rvert^p\), which has finite integral), so, subtracting and applying the second bullet point of 3.43 to the nonnegative integrable functions

\begin{equation*} \bigl(s\lvert u\rvert + (1-s)\lvert v\rvert\bigr)^p - \bigl\lvert su + (1-s)v\bigr\rvert^p \quad\text{and}\quad s\lvert u\rvert^p + (1-s)\lvert v\rvert^p - \bigl(s\lvert u\rvert + (1-s)\lvert v\rvert\bigr)^p , \end{equation*}

each of which has integral \(0\), we conclude that for almost every \(x \in X\):

\begin{equation*} \text{(i)}\quad \bigl\lvert s u(x) + (1-s) v(x)\bigr\rvert = s\lvert u(x)\rvert + (1-s)\lvert v(x)\rvert , \end{equation*}

\begin{equation*} \text{(ii)}\quad \bigl(s \lvert u(x)\rvert + (1-s)\lvert v(x)\rvert\bigr)^p = s\lvert u(x)\rvert^p + (1-s)\lvert v(x)\rvert^p . \end{equation*}

(For (i) we used that \(t \mapsto t^p\) is injective on \([0,\infty)\).)

Fix \(x\) in the set of full measure where both (i) and (ii) hold. By Fact 2, (ii) forces \(\lvert u(x)\rvert = \lvert v(x)\rvert\). By Fact 1 applied to \(z = su(x)\) and \(w = (1-s)v(x)\), condition (i) forces \(s(1-s)\, u(x)\overline{v(x)} \in [0,\infty)\), hence \(u(x)\overline{v(x)} \geq 0\) (a nonnegative real). If \(u(x) = 0\), then \(\lvert v(x)\rvert = 0\) too, so \(u(x) = v(x)\). If \(u(x) \neq 0\), then

\begin{equation*} u(x)\overline{v(x)} = \bigl\lvert u(x)\overline{v(x)}\bigr\rvert = \lvert u(x)\rvert\,\lvert v(x)\rvert = \lvert u(x)\rvert^2 = u(x)\overline{u(x)} , \end{equation*}

and dividing by \(u(x) \neq 0\) gives \(\overline{v(x)} = \overline{u(x)}\), so again \(u(x) = v(x)\).

Thus \(u = v\) almost everywhere, that is, \(f/A = g/B\) almost everywhere, that is,

\begin{equation*} B f(x) = A g(x) \qquad \text{for almost every } x \in X . \end{equation*}

Taking \(a = B > 0\) and \(b = A > 0\) completes the proof.

Problem (7A.21)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(f, g \in \mathcal{L}^1(\mu)\). Prove that

\begin{equation*} \|f + g\|_1 = \|f\|_1 + \|g\|_1 \end{equation*}

if and only if \(f(x)\overline{g(x)} \geq 0\) for almost every \(x \in X\).

Solution

Here the condition \(f(x)\overline{g(x)} \geq 0\) means that the number \(f(x)\overline{g(x)} \in \mathbf{F}\) is a nonnegative real number; when \(\mathbf{F} = \mathbb{R}\) it simply says \(f(x)g(x) \geq 0\).

A pointwise lemma. For \(z, w \in \mathbf{F}\),

\begin{equation*} \lvert z + w \rvert = \lvert z \rvert + \lvert w\rvert \iff z\overline{w} \geq 0 . \end{equation*}

To see this, compute

\begin{equation*} \lvert z + w\rvert^2 = (z+w)\overline{(z+w)} = \lvert z\rvert^2 + 2\operatorname{Re}(z\overline{w}) + \lvert w\rvert^2 \end{equation*}

and

\begin{equation*} \bigl(\lvert z\rvert + \lvert w\rvert\bigr)^2 = \lvert z\rvert^2 + 2\lvert z\rvert\,\lvert w\rvert + \lvert w\rvert^2 = \lvert z\rvert^2 + 2\bigl\lvert z\overline{w}\bigr\rvert + \lvert w\rvert^2 . \end{equation*}

Since \(\lvert z + w\rvert\) and \(\lvert z\rvert + \lvert w\rvert\) are both nonnegative, they are equal if and only if their squares are equal, which happens if and only if

\begin{equation*} \operatorname{Re}\bigl(z\overline{w}\bigr) = \bigl\lvert z\overline{w}\bigr\rvert . \end{equation*}

A complex number \(\zeta\) satisfies \(\operatorname{Re} \zeta = \lvert \zeta\rvert\) if and only if \(\zeta\) is a nonnegative real number (if \(\zeta = \alpha + i\beta\) with \(\alpha, \beta\) real, then \(\alpha = \sqrt{\alpha^2 + \beta^2}\) forces \(\alpha \geq 0\) and \(\beta = 0\); the converse is clear). This proves the lemma.

The main argument. Define \(\psi \colon X \to [0,\infty)\) by

\begin{equation*} \psi(x) = \lvert f(x)\rvert + \lvert g(x) \rvert - \lvert f(x) + g(x)\rvert . \end{equation*}

By the triangle inequality in \(\mathbf{F}\) we have \(\psi \geq 0\) on \(X\). Also \(\psi\) is \(\mathcal{S}\)-measurable, and \(\psi \leq 2(\lvert f\rvert + \lvert g\rvert)\), so \(\psi \in \mathcal{L}^1(\mu)\) because \(f, g \in \mathcal{L}^1(\mu)\). Since \(\lvert f\rvert\), \(\lvert g\rvert\), and \(\lvert f + g\rvert\) each have finite integral, additivity of the integral gives

\begin{equation*} \int \psi \, d\mu = \int \lvert f\rvert \, d\mu + \int \lvert g\rvert\, d\mu - \int \lvert f + g\rvert \, d\mu = \|f\|_1 + \|g\|_1 - \|f + g\|_1 . \end{equation*}

Because \(\psi \geq 0\), we have \(\|\psi\|_1 = \int \psi \, d\mu\). Hence

\begin{equation*} \|f+g\|_1 = \|f\|_1 + \|g\|_1 \iff \|\psi\|_1 = 0 \iff \psi(x) = 0 \ \text{for almost every } x \in X, \end{equation*}

where the last equivalence is the second bullet point of 3.43.

Finally, \(\psi(x) = 0\) means precisely that \(\lvert f(x) + g(x)\rvert = \lvert f(x)\rvert + \lvert g(x)\rvert\), which by the lemma (with \(z = f(x)\) and \(w = g(x)\)) holds if and only if \(f(x)\overline{g(x)} \geq 0\). Therefore

\begin{equation*} \|f + g\|_1 = \|f\|_1 + \|g\|_1 \iff f(x)\overline{g(x)} \geq 0 \ \text{ for almost every } x \in X , \end{equation*}

as desired.

Problem (7A.22)

Suppose \((X, \mathcal{S}, \mu)\) and \((Y, \mathcal{T}, \nu)\) are \(\sigma\)-finite measure spaces and \(0 < p < \infty\). Prove that if \(f \in \mathcal{L}^p(\mu \times \nu)\), then

\begin{equation*} [f]_x \in \mathcal{L}^p(\nu) \quad \text{for almost every } x \in X \end{equation*}

and

\begin{equation*} [f]^y \in \mathcal{L}^p(\mu) \quad \text{for almost every } y \in Y, \end{equation*}

where \([f]_x\) and \([f]^y\) are the cross sections of \(f\) as defined in 5.7.

Solution

Recall that \(f \in \mathcal{L}^p(\mu \times \nu)\) means that \(f\) is an \(\mathcal{S} \otimes \mathcal{T}\)-measurable function on \(X \times Y\) with

\begin{equation*} \|f\|_p^p = \int_{X \times Y} |f|^p \, d(\mu \times \nu) < \infty . \end{equation*}

Measurability of the cross sections. Fix \(x \in X\). Since \(f\) is \(\mathcal{S} \otimes \mathcal{T}\)-measurable, 5.9 tells us that the cross section \([f]_x \colon Y \to \mathbf{F}\) is \(\mathcal{T}\)-measurable; similarly \([f]^y \colon X \to \mathbf{F}\) is \(\mathcal{S}\)-measurable for every \(y \in Y\). Because \(z \mapsto |z|^p\) is a continuous function on \(\mathbf{F}\), the composition \(|[f]_x|^p\) is a \(\mathcal{T}\)-measurable function from \(Y\) to \([0, \infty)\), and \(|[f]^y|^p\) is \(\mathcal{S}\)-measurable. Thus the only thing that needs proof is the finiteness of the relevant integrals for almost every \(x\) and almost every \(y\).

A pointwise identity. The function \(|f|^p \colon X \times Y \to [0, \infty)\) is \(\mathcal{S} \otimes \mathcal{T}\)-measurable (again because \(z \mapsto |z|^p\) is continuous). Directly from the definition 5.7 of cross sections, for all \(x \in X\) and \(y \in Y\) we have

\begin{equation*} \bigl[\,|f|^p\,\bigr]_x(y) = |f(x,y)|^p = \bigl|[f]_x(y)\bigr|^p , \end{equation*}

so \([\,|f|^p\,]_x = |[f]_x|^p\) as functions on \(Y\). Symmetrically, \([\,|f|^p\,]^y = |[f]^y|^p\).

Applying Tonelli’s Theorem. Both measure spaces are \(\sigma\)-finite and \(|f|^p\) is a nonnegative \(\mathcal{S} \otimes \mathcal{T}\)-measurable function on \(X \times Y\). Hence Tonelli’s Theorem (5.28) applies and gives that the function

\begin{equation*} h \colon X \to [0, \infty], \qquad h(x) = \int_Y \bigl|[f]_x\bigr|^p \, d\nu \end{equation*}

is \(\mathcal{S}\)-measurable, and that

\begin{equation*} \int_X h \, d\mu = \int_X \left( \int_Y \bigl[\,|f|^p\,\bigr]_x \, d\nu \right) d\mu = \int_{X \times Y} |f|^p \, d(\mu \times \nu) = \|f\|_p^p < \infty . \end{equation*}

Finiteness almost everywhere. Let

\begin{equation*} E = \{ x \in X : h(x) = \infty \} . \end{equation*}

Since \(h\) is \(\mathcal{S}\)-measurable, \(E \in \mathcal{S}\). For every \(c > 0\) we have \(c \chi_E \leq h\) pointwise on \(X\), so 3.7 (the integral of a simple function) together with 3.8 (integration is order preserving) gives

\begin{equation*} c \, \mu(E) = \int_X c \chi_E \, d\mu \leq \int_X h \, d\mu = \|f\|_p^p . \end{equation*}

Therefore \(\mu(E) \leq \|f\|_p^p / c\) for every \(c > 0\). Because \(\|f\|_p^p < \infty\), letting \(c \to \infty\) yields \(\mu(E) = 0\).

Thus for every \(x \in X \setminus E\) (that is, for almost every \(x \in X\)) the function \([f]_x\) is \(\mathcal{T}\)-measurable and satisfies

\begin{equation*} \int_Y \bigl|[f]_x\bigr|^p \, d\nu = h(x) < \infty , \end{equation*}

which is exactly the statement \([f]_x \in \mathcal{L}^p(\nu)\).

The other cross sections. The argument is symmetric. Part (b) of Tonelli’s Theorem (5.28) says that

\begin{equation*} \tilde{h} \colon Y \to [0, \infty], \qquad \tilde{h}(y) = \int_X \bigl|[f]^y\bigr|^p \, d\mu \end{equation*}

is \(\mathcal{T}\)-measurable with

\begin{equation*} \int_Y \tilde{h} \, d\nu = \int_{X \times Y} |f|^p \, d(\mu \times \nu) = \|f\|_p^p < \infty . \end{equation*}

Repeating the argument of the previous paragraph with \(\tilde{E} = \{ y \in Y : \tilde{h}(y) = \infty \}\) gives \(\nu(\tilde{E}) = 0\), and hence \([f]^y \in \mathcal{L}^p(\mu)\) for almost every \(y \in Y\).

Remark. The \(\sigma\)-finiteness hypothesis is used only to invoke Tonelli’s Theorem, but it is genuinely needed there. Note also that the conclusion cannot be improved to hold for every \(x \in X\). For instance, take \(X = Y = \mathbb{R}\) with \(\mathcal{S} = \mathcal{T}\) the Borel \(\sigma\)-algebra and \(\mu = \nu\) Lebesgue measure, and let

\begin{equation*} f(x, y) = \chi_{\{0\}}(x) \cdot 1 . \end{equation*}

Then \(f = 0\) except on the set \(\{0\} \times \mathbb{R}\), which has product measure \(0\), so \(\|f\|_p = 0\) and \(f \in \mathcal{L}^p(\mu \times \nu)\); yet \([f]_0\) is the constant function \(1\) on \(\mathbb{R}\), which is not in \(\mathcal{L}^p(\nu)\). So the exceptional set of \(x\) in the conclusion can genuinely be nonempty.

Problem (7A.23)

Suppose \(1 \leq p < \infty\) and \(f \in \mathcal{L}^p(\mathbb{R})\).

(a)
For \(t \in \mathbb{R}\), define \(f_t \colon \mathbb{R} \to \mathbb{R}\) by \(f_t(x) = f(x - t)\). Prove that the function \(t \mapsto \|f - f_t\|_p\) is bounded and uniformly continuous on \(\mathbb{R}\).
(b)
For \(t > 0\), define \(f_t \colon \mathbb{R} \to \mathbb{R}\) by \(f_t(x) = f(tx)\). Prove that

\begin{equation*} \lim_{t \to 1} \|f - f_t\|_p = 0 . \end{equation*}

Solution

Throughout, \(| \cdot |\) denotes outer measure on \(\mathbb{R}\) (which agrees with Lebesgue measure on measurable sets), and all integrals are with respect to Lebesgue measure \(\lambda\) on \(\mathbb{R}\). By 7.19 we may take the underlying \(\sigma\)-algebra to be either the Borel sets or the Lebesgue measurable sets; the arguments below work verbatim for either choice, so we simply say measurable. Functions in \(\mathcal{L}^p(\mathbb{R})\) take values in \(\mathbf{F}\); nothing below uses whether \(\mathbf{F}\) is \(\mathbb{R}\) or \(\mathbb{C}\). We first record two invariance properties of Lebesgue measure and the density result that drives both parts.

Lemma 1 (translation invariance of the integral). Suppose \(t \in \mathbb{R}\) and \(h \colon \mathbb{R} \to [0, \infty]\) is measurable. Then \(x \mapsto h(x - t)\) is measurable and

\begin{equation*} \int_{-\infty}^{\infty} h(x - t) \, dx = \int_{-\infty}^{\infty} h . \end{equation*}

Proof. First, translates of measurable sets are measurable, with the same measure. If \(E\) is Borel then \(t + E\) is the inverse image of \(E\) under the continuous (hence Borel measurable) map \(x \mapsto x - t\), so \(t + E\) is Borel. If \(E\) is Lebesgue measurable, then by the definition 2.70 there is a Borel set \(B \subseteq E\) with \(|E \setminus B| = 0\); now \(t + B\) is a Borel subset of \(t + E\) and \(|(t + E) \setminus (t + B)| = |t + (E \setminus B)| = |E \setminus B| = 0\) by 2.7, so \(t + E\) is Lebesgue measurable. In either case \(|t + E| = |E|\) by 2.7, because Lebesgue measure is the restriction of outer measure.

Measurability of \(x \mapsto h(x - t)\) now follows: for every Borel set \(B \subseteq [-\infty, \infty]\) we have \(\{ x \in \mathbb{R} : h(x - t) \in B \} = t + h^{-1}(B)\), which is measurable by the previous paragraph.

For the identity, suppose first that \(h\) is simple, say \(h = \sum_{k=1}^{n} c_k \chi_{E_k}\) with \(E_1, \dots, E_n\) disjoint measurable sets and \(c_1, \dots, c_n \in [0, \infty)\). Since \(\chi_{E_k}(x - t) = \chi_{t + E_k}(x)\), we have \(h(\cdot - t) = \sum_{k=1}^{n} c_k \chi_{t + E_k}\), and the sets \(t + E_1, \dots, t + E_n\) are disjoint and measurable with \(|t + E_k| = |E_k|\). Hence 3.7 gives

\begin{equation*} \int_{-\infty}^{\infty} h(x - t) \, dx = \sum_{k=1}^{n} c_k \, |t + E_k| = \sum_{k=1}^{n} c_k \, |E_k| = \int_{-\infty}^{\infty} h . \end{equation*}

For general measurable \(h \geq 0\), use 2.89 to choose simple measurable functions \(h_1, h_2, \dots\) with \(0 \leq h_1 \leq h_2 \leq \cdots \leq h\) and \(h_k \to h\) pointwise. Then \(h_k(\cdot - t)\) is an increasing sequence of nonnegative simple measurable functions converging pointwise to \(h(\cdot - t)\), so two applications of the Monotone Convergence Theorem (3.11), together with the identity for simple functions, give

\begin{equation*} \int_{-\infty}^{\infty} h(x - t) \, dx = \lim_{k \to \infty} \int_{-\infty}^{\infty} h_k(x - t) \, dx = \lim_{k \to \infty} \int_{-\infty}^{\infty} h_k = \int_{-\infty}^{\infty} h . \qquad \square \end{equation*}

Lemma 2 (dilation). Suppose \(t > 0\) and \(h \colon \mathbb{R} \to [0, \infty]\) is measurable. Then \(x \mapsto h(tx)\) is measurable and

\begin{equation*} \int_{-\infty}^{\infty} h(tx) \, dx = \frac{1}{t} \int_{-\infty}^{\infty} h . \end{equation*}

Proof. For \(c > 0\) and \(A \subseteq \mathbb{R}\) write \(cA = \{ ca : a \in A \}\). The map \(I \mapsto cI\) is a bijection of the set of open intervals onto itself with \(\ell(cI) = c \, \ell(I)\), and \(A \subseteq \bigcup_k I_k\) if and only if \(cA \subseteq \bigcup_k cI_k\); hence the definition 2.2 of outer measure gives

\begin{equation*} |cA| = c \, |A| \qquad \text{for every } A \subseteq \mathbb{R} \text{ and every } c > 0 . \end{equation*}

Next, dilates of measurable sets are measurable. If \(E\) is Borel then \(cE\) is the inverse image of \(E\) under the continuous map \(x \mapsto x / c\), so \(cE\) is Borel. If \(E\) is Lebesgue measurable, choose by 2.70 a Borel set \(B \subseteq E\) with \(|E \setminus B| = 0\); then \(cB\) is a Borel subset of \(cE\) and \(|(cE) \setminus (cB)| = |c(E \setminus B)| = c \, |E \setminus B| = 0\), so \(cE\) is Lebesgue measurable. In either case \(|cE| = c|E|\).

Measurability of \(x \mapsto h(tx)\) follows because \(\{ x \in \mathbb{R} : h(tx) \in B \} = t^{-1} h^{-1}(B)\) for every Borel set \(B \subseteq [-\infty, \infty]\).

For the identity, note that \(\chi_E(tx) = \chi_{t^{-1}E}(x)\) and \(|t^{-1}E| = t^{-1}|E|\). Hence if \(h = \sum_{k=1}^{n} c_k \chi_{E_k}\) with \(E_1, \dots, E_n\) disjoint measurable sets and \(c_1, \dots, c_n \in [0, \infty)\), then \(h(t \, \cdot) = \sum_{k=1}^{n} c_k \chi_{t^{-1}E_k}\) with \(t^{-1}E_1, \dots, t^{-1}E_n\) disjoint and measurable, so 3.7 gives

\begin{equation*} \int_{-\infty}^{\infty} h(tx) \, dx = \sum_{k=1}^{n} c_k \, t^{-1} |E_k| = \frac{1}{t} \int_{-\infty}^{\infty} h . \end{equation*}

The extension to all measurable \(h \geq 0\) is by 2.89 and two applications of the Monotone Convergence Theorem (3.11), exactly as in Lemma 1. \(\square\)

The measurability arguments in the two proofs above apply verbatim to \(\mathbf{F}\)-valued functions: if \(u \colon \mathbb{R} \to \mathbf{F}\) is measurable, then so are \(x \mapsto u(x - t)\) and \(x \mapsto u(tx)\), since their inverse images of a Borel subset of \(\mathbf{F}\) are \(t + u^{-1}(B)\) and \(t^{-1}u^{-1}(B)\) respectively. In particular the functions \(f_t\) of parts (a) and (b) are measurable.

Applying Lemma 1 to \(h = |f|^p\) shows that in part (a) we have \(f_t \in \mathcal{L}^p(\mathbb{R})\) with

\begin{equation*} \|f_t\|_p = \|f\|_p \qquad \text{for all } t \in \mathbb{R}, \end{equation*}

and applying Lemma 2 to \(h = |f|^p\) shows that in part (b) we have \(f_t \in \mathcal{L}^p(\mathbb{R})\) with

\begin{equation*} \|f_t\|_p = t^{-1/p} \|f\|_p \qquad \text{for all } t > 0 . \end{equation*}

Lemma 3 (density). For every \(\varepsilon > 0\) there is a continuous function \(g \colon \mathbb{R} \to \mathbf{F}\) such that \(\|f - g\|_p < \varepsilon\) and \(\{ x \in \mathbb{R} : g(x) \neq 0 \}\) is bounded. This is precisely Exercise 19 of this section (which extends 3.48 from \(\mathcal{L}^1(\mathbb{R})\) to \(\mathcal{L}^p(\mathbb{R})\)). Note that such a \(g\) is bounded and vanishes outside some interval \([-M, M]\), so \(g \in \mathcal{L}^p(\mathbb{R})\).

Lemma 4 (such \(g\) is uniformly continuous). Suppose \(g \colon \mathbb{R} \to \mathbf{F}\) is continuous and \(g(x) = 0\) for all \(x\) with \(|x| > M\), where \(M > 0\). Then \(g\) is uniformly continuous on \(\mathbb{R}\).

Proof. Let \(\eta > 0\). Since \(g\) is continuous on the compact interval \([-M-2, M+2]\), it is uniformly continuous there, so there exists \(\delta \in (0, 1)\) such that \(|g(a) - g(b)| < \eta\) whenever \(a, b \in [-M-2, M+2]\) and \(|a - b| < \delta\). Now suppose \(a, b \in \mathbb{R}\) with \(|a - b| < \delta\). If both lie in \([-M-2, M+2]\) we are done. Otherwise one of them, say \(a\), satisfies \(|a| > M + 2\); then \(|b| > M + 1 > M\), so \(g(a) = g(b) = 0\) and again \(|g(a) - g(b)| < \eta\). \(\square\)

Part (a). Define \(\varphi \colon \mathbb{R} \to [0, \infty)\) by \(\varphi(t) = \|f - f_t\|_p\).

Boundedness. By Minkowski’s inequality (7.14) and translation invariance,

\begin{equation*} \varphi(t) \leq \|f\|_p + \|f_t\|_p = 2 \|f\|_p < \infty \end{equation*}

for every \(t \in \mathbb{R}\). Thus \(\varphi\) is bounded by \(2\|f\|_p\).

A reduction to continuity at \(0\). For \(a, b \in \mathbb{R}\) and \(x \in \mathbb{R}\) we have \((f_a)_b(x) = f_a(x - b) = f(x - a - b)\), so \((f_a)_b = f_{a+b}\). Hence, for \(s, t \in \mathbb{R}\),

\begin{equation*} (f - f_{s-t})_t(x) = f(x - t) - f\bigl(x - t - (s-t)\bigr) = f_t(x) - f_s(x), \end{equation*}

so \(f_t - f_s = (f - f_{s-t})_t\) and therefore, by translation invariance,

\begin{equation*} \|f_t - f_s\|_p = \|f - f_{s-t}\|_p = \varphi(s - t) . \end{equation*}

Minkowski’s inequality (7.14) applied twice gives \(\bigl| \|f - f_t\|_p - \|f - f_s\|_p \bigr| \leq \|f_t - f_s\|_p\), and thus

\begin{equation*} |\varphi(t) - \varphi(s)| \leq \varphi(s - t) \qquad \text{for all } s, t \in \mathbb{R} . \end{equation*}

Since \(\varphi(0) = 0\), it now suffices to prove that \(\lim_{u \to 0} \varphi(u) = 0\): given \(\varepsilon > 0\), a \(\delta > 0\) with \(\varphi(u) < \varepsilon\) for all \(|u| < \delta\) immediately yields \(|\varphi(t) - \varphi(s)| < \varepsilon\) whenever \(|s - t| < \delta\), which is uniform continuity of \(\varphi\) on \(\mathbb{R}\).

Proof that \(\lim_{u \to 0} \varphi(u) = 0\). Let \(\varepsilon > 0\). By Lemma 3 choose a continuous \(g \colon \mathbb{R} \to \mathbf{F}\) vanishing outside \([-M, M]\) for some \(M > 0\) and satisfying \(\|f - g\|_p < \varepsilon / 3\). For \(u \in \mathbb{R}\) we have \((f - g)_u = f_u - g_u\), so translation invariance gives \(\|f_u - g_u\|_p = \|f - g\|_p < \varepsilon / 3\). Minkowski’s inequality (7.14) then gives

\begin{equation*} \varphi(u) = \|f - f_u\|_p \leq \|f - g\|_p + \|g - g_u\|_p + \|g_u - f_u\|_p < \tfrac{2\varepsilon}{3} + \|g - g_u\|_p . \end{equation*}

Now suppose \(|u| \leq 1\). If \(|x| > M + 1\) then \(|x| > M\) and \(|x - u| > M\), so \(g(x) - g_u(x) = 0\); thus \(g - g_u\) vanishes outside \([-M-1, M+1]\) and hence

\begin{equation*} \|g - g_u\|_p^p = \int_{-M-1}^{M+1} |g(x) - g(x - u)|^p \, dx \leq (2M + 2) \Bigl( \sup_{x \in \mathbb{R}} |g(x) - g(x-u)| \Bigr)^p . \end{equation*}

By Lemma 4 the supremum on the right tends to \(0\) as \(u \to 0\). Hence there exists \(\delta \in (0, 1]\) such that \(\|g - g_u\|_p < \varepsilon / 3\) whenever \(|u| < \delta\), and then \(\varphi(u) < \varepsilon\). This proves \(\lim_{u \to 0} \varphi(u) = 0\) and completes part (a).

Part (b). Now \(f_t(x) = f(tx)\) for \(t > 0\). Let \(\varepsilon > 0\). By Lemma 3 choose a continuous \(g \colon \mathbb{R} \to \mathbf{F}\) vanishing outside \([-M, M]\) for some \(M > 0\) with \(\|f - g\|_p < \varepsilon\). Writing \(g_t(x) = g(tx)\), we have \((f - g)_t = f_t - g_t\), so Lemma 2 gives

\begin{equation*} \|f_t - g_t\|_p = t^{-1/p} \|f - g\|_p < t^{-1/p} \varepsilon \leq 2^{1/p} \varepsilon \leq 2 \varepsilon \end{equation*}

for all \(t \in [\tfrac12, 2]\) (using \(p \geq 1\)). Minkowski’s inequality (7.14) then gives, for such \(t\),

\begin{equation*} \|f - f_t\|_p \leq \|f - g\|_p + \|g - g_t\|_p + \|g_t - f_t\|_p < 3\varepsilon + \|g - g_t\|_p . \end{equation*}

It remains to show that \(\|g - g_t\|_p \to 0\) as \(t \to 1\). Suppose \(t \in [\tfrac12, 2]\). If \(|x| > 2M\) then \(|x| > M\) and \(|tx| \geq |x| / 2 > M\), so \(g(x) - g(tx) = 0\); hence \(g - g_t\) vanishes outside \([-2M, 2M]\) and

\begin{equation*} \|g - g_t\|_p^p = \int_{-2M}^{2M} |g(x) - g(tx)|^p \, dx \leq 4M \Bigl( \sup_{|x| \leq 2M} |g(x) - g(tx)| \Bigr)^p . \end{equation*}

Let \(\eta > 0\). By Lemma 4 there exists \(\delta_0 > 0\) such that \(|g(a) - g(b)| < \eta\) whenever \(|a - b| < \delta_0\). If \(|x| \leq 2M\) and \(|t - 1| < \delta_0 / (2M)\), then \(|tx - x| = |t - 1| \, |x| < \delta_0\), so \(|g(x) - g(tx)| < \eta\). Hence

\begin{equation*} \sup_{|x| \leq 2M} |g(x) - g(tx)| \longrightarrow 0 \quad \text{as } t \to 1, \end{equation*}

and therefore \(\|g - g_t\|_p \to 0\) as \(t \to 1\).

Combining the last two paragraphs: choose \(\delta \in (0, \tfrac12)\) so small that \(\|g - g_t\|_p < \varepsilon\) whenever \(|t - 1| < \delta\). Then \(|t - 1| < \delta\) implies \(t \in [\tfrac12, 2]\) and

\begin{equation*} \|f - f_t\|_p < 4 \varepsilon . \end{equation*}

Since \(\varepsilon > 0\) was arbitrary, \(\lim_{t \to 1} \|f - f_t\|_p = 0\).

Problem (7A.24)

Suppose \(1 \leq p < \infty\) and \(f \in \mathcal{L}^p(\mathbb{R})\). Prove that

\begin{equation*} \lim_{t \downarrow 0} \frac{1}{2t} \int_{b-t}^{b+t} |f - f(b)|^p = 0 \end{equation*}

for almost every \(b \in \mathbb{R}\).

Solution

For \(p = 1\) this is precisely the Lebesgue Differentiation Theorem, first version (4.10). The point of the exercise is to handle a general exponent \(p \in [1, \infty)\), and the proof below reduces everything to 4.10 by applying that theorem to the countably many functions \(|f - r|^p\) (suitably localized), where \(r\) ranges over a countable dense subset of \(\mathbf{F}\).

An elementary inequality. If \(p \geq 1\) then \(s \mapsto s^p\) is convex on \([0, \infty)\), so for all \(u, v \in \mathbf{F}\),

\begin{equation*} |u + v|^p \leq \bigl( |u| + |v| \bigr)^p = 2^p \left( \frac{|u| + |v|}{2} \right)^{\!p} \leq 2^p \cdot \frac{|u|^p + |v|^p}{2} = 2^{p-1} \bigl( |u|^p + |v|^p \bigr) . \end{equation*}

A countable dense set. Let \(D\) be a countable dense subset of \(\mathbf{F}\): take \(D = \mathbb{Q}\) if \(\mathbf{F} = \mathbb{R}\), and \(D = \{ q + i q^{\prime} : q, q^{\prime} \in \mathbb{Q} \}\) if \(\mathbf{F} = \mathbb{C}\).

Step 1: localized truncations lie in \(\mathcal{L}^1(\mathbb{R})\). For \(r \in D\) and \(n \in \mathbb{Z}^+\) define \(h_{r,n} \colon \mathbb{R} \to [0, \infty)\) by

\begin{equation*} h_{r,n}(x) = |f(x) - r|^p \, \chi_{(-n, n)}(x) . \end{equation*}

Each \(h_{r,n}\) is Lebesgue measurable, and by the elementary inequality above,

\begin{equation*} \int_{-\infty}^{\infty} h_{r,n} = \int_{-n}^{n} |f - r|^p \leq 2^{p-1} \left( \int_{-n}^{n} |f|^p + 2n |r|^p \right) \leq 2^{p-1} \bigl( \|f\|_p^p + 2n |r|^p \bigr) < \infty . \end{equation*}

Thus \(h_{r,n} \in \mathcal{L}^1(\mathbb{R})\).

Step 2: apply the Lebesgue Differentiation Theorem. By 4.10 applied to \(h_{r,n}\), there is a set \(N_{r,n} \subseteq \mathbb{R}\) with \(|N_{r,n}| = 0\) such that

\begin{equation*} \lim_{t \downarrow 0} \frac{1}{2t} \int_{b-t}^{b+t} \bigl| h_{r,n} - h_{r,n}(b) \bigr| = 0 \qquad \text{for every } b \in \mathbb{R} \setminus N_{r,n} . \end{equation*}

Let

\begin{equation*} N = \bigcup_{r \in D} \ \bigcup_{n \in \mathbb{Z}^+} N_{r,n} . \end{equation*}

This is a countable union of sets of measure \(0\), so \(|N| = 0\) by countable subadditivity of outer measure (2.8).

Step 3: averages of \(|f - r|^p\) converge, for every \(b \notin N\) and every \(r \in D\). Fix \(b \in \mathbb{R} \setminus N\) and \(r \in D\). Choose \(n \in \mathbb{Z}^+\) with \(n > |b| + 1\). Then for every \(t \in (0, 1)\) we have \((b - t, b + t) \subseteq (-n, n)\), so \(h_{r,n}\) agrees with \(|f - r|^p\) on \((b-t, b+t)\); also \(b \in (-n, n)\), so \(h_{r,n}(b) = |f(b) - r|^p\). Since \(\frac{1}{2t} \int_{b-t}^{b+t} |f(b) - r|^p = |f(b) - r|^p\), we obtain for \(0 < t < 1\)

\begin{equation*} \left| \frac{1}{2t} \int_{b-t}^{b+t} |f - r|^p - |f(b) - r|^p \right| \leq \frac{1}{2t} \int_{b-t}^{b+t} \Bigl| \, |f - r|^p - |f(b) - r|^p \Bigr| = \frac{1}{2t} \int_{b-t}^{b+t} \bigl| h_{r,n} - h_{r,n}(b) \bigr| , \end{equation*}

and the right side tends to \(0\) as \(t \downarrow 0\) because \(b \notin N_{r,n}\). Hence

\begin{equation*} \lim_{t \downarrow 0} \frac{1}{2t} \int_{b-t}^{b+t} |f - r|^p = |f(b) - r|^p \qquad \text{for every } r \in D . \end{equation*}

Step 4: from the dense set to \(f(b)\). Continue to fix \(b \in \mathbb{R} \setminus N\), and let \(\varepsilon > 0\). Since \(f(b) \in \mathbf{F}\) and \(D\) is dense in \(\mathbf{F}\), there exists \(r \in D\) with \(|f(b) - r| < \varepsilon\). For every \(x \in \mathbb{R}\), the elementary inequality applied to \(u = f(x) - r\) and \(v = r - f(b)\) gives

\begin{equation*} |f(x) - f(b)|^p \leq 2^{p-1} \bigl( |f(x) - r|^p + |r - f(b)|^p \bigr) . \end{equation*}

Integrating over \((b - t, b + t)\) and dividing by \(2t\) yields, for every \(t \in (0, 1)\),

\begin{equation*} \frac{1}{2t} \int_{b-t}^{b+t} |f - f(b)|^p \leq 2^{p-1} \left( \frac{1}{2t} \int_{b-t}^{b+t} |f - r|^p + |r - f(b)|^p \right) . \end{equation*}

(All these integrals are finite: the right side is finite by Step 1, hence so is the left side.) Letting \(t \downarrow 0\) and using Step 3,

\begin{equation*} \limsup_{t \downarrow 0} \ \frac{1}{2t} \int_{b-t}^{b+t} |f - f(b)|^p \leq 2^{p-1} \bigl( |f(b) - r|^p + |r - f(b)|^p \bigr) = 2^p |f(b) - r|^p < 2^p \varepsilon^p . \end{equation*}

Because \(\varepsilon > 0\) was arbitrary and the quantity on the left does not depend on \(\varepsilon\), the limsup equals \(0\). As the averages are nonnegative, this means

\begin{equation*} \lim_{t \downarrow 0} \frac{1}{2t} \int_{b-t}^{b+t} |f - f(b)|^p = 0 . \end{equation*}

Since this holds for every \(b \in \mathbb{R} \setminus N\) and \(|N| = 0\), the conclusion holds for almost every \(b \in \mathbb{R}\), as desired.

Remark. The only place the hypothesis \(f \in \mathcal{L}^p(\mathbb{R})\) was used is in Step 1, to guarantee that \(|f - r|^p\) is integrable on each bounded interval. So the same proof shows the conclusion for every Lebesgue measurable \(f\) with \(\int_{-n}^n |f|^p < \infty\) for all \(n \in \mathbb{Z}^+\). A point \(b\) at which the displayed limit is \(0\) is called a Lebesgue point of \(f\) for the exponent \(p\). The essential trick is that a single null set \(N\) works for all \(r \in D\) at once: it is the countability of \(D\) that allows the exceptional sets \(N_{r,n}\) to be assembled into one null set before the density of \(D\) is used to replace \(r\) by \(f(b)\).

Exercises 7B

Problem (7B.1)

Suppose \(n > 1\) and \(0 < p < 1\). Prove that if \(\|\cdot\|\) is defined on \(\mathbb{F}^n\) by

\begin{equation*} \|(a_1, \ldots, a_n)\| = \bigl(|a_1|^p + \cdots + |a_n|^p\bigr)^{1/p}, \end{equation*}

then \(\|\cdot\|\) is not a norm on \(\mathbb{F}^n\).

Solution

We show that the triangle inequality fails; because \(n > 1\) there is room in \(\mathbb{F}^n\) for two vectors with disjoint supports, and because \(p < 1\) the exponent \(1/p\) is larger than \(1\).

Let

\begin{equation*} a = (1, 0, 0, \ldots, 0), \qquad b = (0, 1, 0, \ldots, 0), \end{equation*}

both of which lie in \(\mathbb{F}^n\) because \(n \ge 2\). Then

\begin{equation*} \|a\| = (1^p)^{1/p} = 1 \qquad\text{and}\qquad \|b\| = (1^p)^{1/p} = 1, \end{equation*}

so \(\|a\| + \|b\| = 2\). On the other hand \(a + b = (1, 1, 0, \ldots, 0)\), so

\begin{equation*} \|a + b\| = (1^p + 1^p)^{1/p} = 2^{1/p}. \end{equation*}

Because \(0 < p < 1\) we have \(1/p > 1\), and hence \(2^{1/p} > 2^1 = 2\). Therefore

\begin{equation*} \|a + b\| = 2^{1/p} > 2 = \|a\| + \|b\|, \end{equation*}

so the triangle inequality fails for the pair \(a, b\). A norm must satisfy the triangle inequality, so \(\|\cdot\|\) is not a norm on \(\mathbb{F}^n\).

Remark: the failure is exactly the triangle inequality, and nothing else. Indeed \(\|c\| = 0\) if and only if \(c = 0\), and homogeneity holds, since for \(\alpha \in \mathbb{F}\)

\begin{equation*} \|\alpha(a_1, \ldots, a_n)\| = \bigl(|\alpha|^p |a_1|^p + \cdots + |\alpha|^p |a_n|^p\bigr)^{1/p} = |\alpha| \, \|(a_1, \ldots, a_n)\|. \end{equation*}

Note also that both hypotheses are needed: if \(n = 1\) then \(\|(a_1)\| = |a_1|\), which is a norm; and if \(p \ge 1\) then \(\|\cdot\|\) is a norm on \(\mathbb{F}^n\) by Minkowski’s inequality (7.14) applied to counting measure on \(\{1, \ldots, n\}\).

Problem (7B.2)

(a) Suppose \(1 \le p < \infty\). Prove that there is a countable subset of \(\ell^p\) whose closure equals \(\ell^p\).

(b) Prove that there does not exist a countable subset of \(\ell^\infty\) whose closure equals \(\ell^\infty\).

Solution

Throughout, let \(\mathbb{D}\) denote \(\mathbb{Q}\) if \(\mathbb{F} = \mathbb{R}\) and \(\mathbb{Q} + i\mathbb{Q}\) if \(\mathbb{F} = \mathbb{C}\). In either case \(\mathbb{D}\) is a countable dense subset of \(\mathbb{F}\). Recall (see the discussion following 7.18) that \(\ell^p\) is \(L^p(\mu) = \mathcal{L}^p(\mu)\) for \(\mu\) counting measure on \(\mathbb{Z}^+\); because counting measure has no nonempty sets of measure \(0\), no identification of functions occurs and the elements of \(\ell^p\) are genuine sequences.

(a) Let

\begin{equation*} D = \{(a_1, a_2, \ldots) : \text{each } a_k \in \mathbb{D} \text{ and } a_k = 0 \text{ for all but finitely many } k\}. \end{equation*}

For each \(n \in \mathbb{Z}^+\), the set \(D_n\) of elements of \(D\) that vanish after the \(n^{\text{th}}\) coordinate is in bijection with \(\mathbb{D}^n\), hence is countable; and \(D = \bigcup_{n=1}^{\infty} D_n\) is a countable union of countable sets, hence countable. Every element of \(D\) has only finitely many nonzero terms, so \(D \subseteq \ell^p\).

Now suppose \(a = (a_1, a_2, \ldots) \in \ell^p\) and \(\varepsilon > 0\). Because \(\sum_{k=1}^{\infty} |a_k|^p < \infty\), the tail of this convergent series tends to \(0\); thus there exists \(n \in \mathbb{Z}^+\) such that

\begin{equation*} \sum_{k=n+1}^{\infty} |a_k|^p < \frac{\varepsilon^p}{2}. \end{equation*}

For each \(k \in \{1, \ldots, n\}\), use the density of \(\mathbb{D}\) in \(\mathbb{F}\) to choose \(b_k \in \mathbb{D}\) with

\begin{equation*} |a_k - b_k|^p < \frac{\varepsilon^p}{2n}. \end{equation*}

Let \(b = (b_1, \ldots, b_n, 0, 0, \ldots) \in D\). Then

\begin{equation*} \|a - b\|_p^p = \sum_{k=1}^{n} |a_k - b_k|^p + \sum_{k=n+1}^{\infty} |a_k|^p < n \cdot \frac{\varepsilon^p}{2n} + \frac{\varepsilon^p}{2} = \varepsilon^p, \end{equation*}

so \(\|a - b\|_p < \varepsilon\). Thus every element of \(\ell^p\) is in the closure of \(D\); since \(D \subseteq \ell^p\) and \(\ell^p\) is closed in itself, the closure of \(D\) equals \(\ell^p\).

(b) Let

\begin{equation*} S = \{(a_1, a_2, \ldots) : a_k \in \{0, 1\} \text{ for every } k \in \mathbb{Z}^+\}. \end{equation*}

Each element of \(S\) is a bounded sequence, so \(S \subseteq \ell^\infty\). The set \(S\) is uncountable: it is in bijection with the set of all subsets of \(\mathbb{Z}^+\) (a sequence of \(0\)s and \(1\)s corresponds to the set of indices where it equals \(1\)), and the set of all subsets of \(\mathbb{Z}^+\) is uncountable by Cantor’s theorem.

If \(a, b \in S\) with \(a \ne b\), then \(a_k \ne b_k\) for some \(k\), and then \(|a_k - b_k| = 1\) while \(|a_j - b_j| \le 1\) for every \(j\); hence

\begin{equation*} \|a - b\|_\infty = \sup_{j \in \mathbb{Z}^+} |a_j - b_j| = 1. \end{equation*}

So the elements of \(S\) are pairwise at distance exactly \(1\).

Now suppose \(C \subseteq \ell^\infty\) is a set whose closure equals \(\ell^\infty\). For each \(a \in S\) we have \(a \in \ell^\infty\), which is the closure of \(C\), so there exists \(c(a) \in C\) with \(\|a - c(a)\|_\infty < \frac{1}{2}\). The map \(a \mapsto c(a)\) from \(S\) to \(C\) is injective: if \(a, b \in S\) and \(c(a) = c(b)\), then

\begin{equation*} \|a - b\|_\infty \le \|a - c(a)\|_\infty + \|c(b) - b\|_\infty < \tfrac{1}{2} + \tfrac{1}{2} = 1, \end{equation*}

which forces \(a = b\) because distinct elements of \(S\) are at distance \(1\). Thus there is an injection from the uncountable set \(S\) into \(C\), so \(C\) is uncountable.

Hence no countable subset of \(\ell^\infty\) has closure equal to \(\ell^\infty\).

Problem (7B.3)

(a) Suppose \(1 \le p < \infty\). Prove that there is a countable subset of \(L^p(\mathbb{R})\) whose closure equals \(L^p(\mathbb{R})\).

(b) Prove that there does not exist a countable subset of \(L^\infty(\mathbb{R})\) whose closure equals \(L^\infty(\mathbb{R})\).

Solution

As in the previous exercise, let \(\mathbb{D} = \mathbb{Q}\) if \(\mathbb{F} = \mathbb{R}\) and \(\mathbb{D} = \mathbb{Q} + i\mathbb{Q}\) if \(\mathbb{F} = \mathbb{C}\); in either case \(\mathbb{D}\) is a countable dense subset of \(\mathbb{F}\). Write \(\lambda\) for Lebesgue measure. To lighten notation we follow the usual practice of denoting an element of \(L^p(\mathbb{R})\) by a function representing it; two functions that agree almost everywhere represent the same element and have the same \(p\)-norm.

(a) Let

\begin{equation*} D = \Bigl\{\sum_{j=1}^{m} c_j \chi_{(a_j, b_j)} : m \in \mathbb{Z}^+,\ c_j \in \mathbb{D},\ a_j, b_j \in \mathbb{Q},\ a_j < b_j\Bigr\}. \end{equation*}

For each fixed \(m\) the set of such sums is the image of the countable set \((\mathbb{D} \times \mathbb{Q} \times \mathbb{Q})^m\) under a map, hence countable; so \(D\) is countable, being a countable union of countable sets. Each element of \(D\) is a bounded measurable function vanishing outside a bounded set, so \(D \subseteq \mathcal{L}^p(\mathbb{R})\). We show that the closure of \(D\) in \(L^p(\mathbb{R})\) equals \(L^p(\mathbb{R})\), in three steps.

Step 1: simple functions are dense. Suppose \(f \in \mathcal{L}^p(\mathbb{R})\). First suppose \(\mathbb{F} = \mathbb{R}\). By 2.89 there is a sequence \(f_1, f_2, \ldots\) of simple Borel (or Lebesgue) measurable functions with \(|f_k(x)| \le |f(x)|\) and \(\lim_{k \to \infty} f_k(x) = f(x)\) for every \(x\). Then

\begin{equation*} |f - f_k|^p \le (|f| + |f_k|)^p \le 2^p |f|^p, \end{equation*}

and \(2^p |f|^p\) is integrable because \(f \in \mathcal{L}^p\). Since \(|f - f_k|^p \to 0\) pointwise, the Dominated Convergence Theorem (3.31) gives \(\int |f - f_k|^p \, d\lambda \to 0\), that is, \(\|f - f_k\|_p \to 0\). If \(\mathbb{F} = \mathbb{C}\), apply the above to \(\operatorname{Re} f\) and \(\operatorname{Im} f\) (each of which lies in \(\mathcal{L}^p\) because \(|\operatorname{Re} f|, |\operatorname{Im} f| \le |f|\)) and add, using Minkowski’s inequality (7.14). In either case, simple functions belonging to \(\mathcal{L}^p(\mathbb{R})\) are dense in \(L^p(\mathbb{R})\).

Step 2: characteristic functions of finite-measure sets are approximated by elements of \(D\). Suppose \(E \subseteq \mathbb{R}\) is measurable with \(\lambda(E) < \infty\), and suppose \(\varepsilon > 0\). By 2.71(e) there is an open set \(G \supseteq E\) with \(\lambda(G \setminus E) < \varepsilon\). Every open subset of \(\mathbb{R}\) is the union of a countable disjoint collection of open intervals (namely its connected components: each component of an open subset of \(\mathbb{R}\) is an open interval, and the components are disjoint and each contains a rational number, so there are at most countably many). Write \(G = \bigcup_{k} J_k\) with \(J_1, J_2, \ldots\) disjoint open intervals. By countable additivity,

\begin{equation*} \sum_{k} \lambda(J_k) = \lambda(G) \le \lambda(E) + \lambda(G \setminus E) < \infty, \end{equation*}

so each \(J_k\) is a bounded interval and there exists \(m\) with \(\sum_{k > m} \lambda(J_k) < \varepsilon\). Put \(U = J_1 \cup \cdots \cup J_m\), so that \(\chi_U = \sum_{k=1}^{m} \chi_{J_k}\) because the \(J_k\) are disjoint. Now

\begin{equation*} E \setminus U \subseteq G \setminus U = \bigcup_{k > m} J_k \quad\text{and}\quad U \setminus E \subseteq G \setminus E, \end{equation*}

so \(\lambda(E \bigtriangleup U) < 2\varepsilon\) and hence

\begin{equation*} \|\chi_E - \chi_U\|_p^p = \int |\chi_E - \chi_U|^p \, d\lambda = \lambda(E \bigtriangleup U) < 2\varepsilon. \end{equation*}

Finally, fix \(\delta > 0\), write \(J_k = (a_k, b_k)\) with \(-\infty < a_k < b_k < \infty\), and choose rationals

\begin{equation*} a_k^{\prime} \in (a_k - \delta, a_k) \qquad\text{and}\qquad b_k^{\prime} \in (b_k, b_k + \delta), \end{equation*}

which is possible because \(\mathbb{Q}\) is dense in \(\mathbb{R}\). Then \(a_k^{\prime} < a_k < b_k < b_k^{\prime}\), so in particular \(a_k^{\prime} < b_k^{\prime}\) and \((a_k^{\prime}, b_k^{\prime}) \supseteq J_k\), whence

\begin{equation*} J_k \bigtriangleup (a_k^{\prime}, b_k^{\prime}) = (a_k^{\prime}, b_k^{\prime}) \setminus J_k = (a_k^{\prime}, a_k] \cup [b_k, b_k^{\prime}) \end{equation*}

has measure less than \(2\delta\); thus \(\|\chi_{J_k} - \chi_{(a_k^{\prime}, b_k^{\prime})}\|_p^p \le 2\delta\). Taking \(u = \sum_{k=1}^{m} \chi_{(a_k^{\prime}, b_k^{\prime})} \in D\) and using Minkowski’s inequality (7.14),

\begin{equation*} \|\chi_U - u\|_p \le \sum_{k=1}^{m} \|\chi_{J_k} - \chi_{(a_k^{\prime}, b_k^{\prime})}\|_p \le m (2\delta)^{1/p}, \end{equation*}

which is less than \(\varepsilon\) for \(\delta\) small. Hence \(\|\chi_E - u\|_p < (2\varepsilon)^{1/p} + \varepsilon\), which can be made arbitrarily small. So \(\chi_E\) lies in the closure of \(D\) whenever \(\lambda(E) < \infty\).

Step 3: putting it together. Suppose \(g\) is a simple function in \(\mathcal{L}^p(\mathbb{R})\), say \(g = \sum_{j=1}^{n} c_j \chi_{E_j}\) where \(c_1, \ldots, c_n\) are the distinct nonzero values of \(g\) and \(E_j = g^{-1}(\{c_j\})\) (see the discussion preceding 2.89). Because \(\int |g|^p \, d\lambda = \sum_{j} |c_j|^p \lambda(E_j) < \infty\) and each \(c_j \ne 0\), we have \(\lambda(E_j) < \infty\) for each \(j\). Given \(\varepsilon > 0\), use Step 2 to choose \(u_j \in D\) with \(\|\chi_{E_j} - u_j\|_p < \varepsilon\), and choose \(q_j \in \mathbb{D}\) with \(|c_j - q_j| < \varepsilon\). Then \(\sum_{j=1}^{n} q_j u_j \in D\) (since \(D\) is closed under sums and under multiplication by elements of \(\mathbb{D}\)), and by Minkowski’s inequality (7.14),

\begin{equation*} \Bigl\| g - \sum_{j=1}^{n} q_j u_j \Bigr\|_p \le \sum_{j=1}^{n} |c_j| \, \|\chi_{E_j} - u_j\|_p + \sum_{j=1}^{n} |c_j - q_j| \, \|u_j\|_p \le \varepsilon \sum_{j=1}^{n} |c_j| + \varepsilon \sum_{j=1}^{n} \|u_j\|_p . \end{equation*}

If \(\varepsilon \le 1\) then \(\|u_j\|_p \le \|\chi_{E_j}\|_p + \varepsilon \le \|\chi_{E_j}\|_p + 1\), so the right side above is at most \(\varepsilon M\), where

\begin{equation*} M = \sum_{j=1}^{n} |c_j| + \sum_{j=1}^{n} \bigl( \|\chi_{E_j}\|_p + 1 \bigr) \end{equation*}

is a finite constant depending only on \(g\) (finite because each \(\lambda(E_j) < \infty\)). Letting \(\varepsilon \to 0\) shows that \(g\) lies in the closure of \(D\).

Combining Steps 1 and 3: given \(f \in \mathcal{L}^p(\mathbb{R})\) and \(\varepsilon > 0\), Step 1 gives a simple \(g \in \mathcal{L}^p(\mathbb{R})\) with \(\|f - g\|_p < \varepsilon/2\) and Step 3 gives \(u \in D\) with \(\|g - u\|_p < \varepsilon/2\); thus \(\|f - u\|_p < \varepsilon\). Therefore the closure of \(D\) (more precisely, of \(\{\widetilde{u} : u \in D\}\)) equals \(L^p(\mathbb{R})\), and \(D\) is countable.

(b) For \(t \in (0, 1)\) let \(g_t = \chi_{(0, t)}\), and let \(S = \{g_t : t \in (0,1)\} \subseteq L^\infty(\mathbb{R})\). If \(0 < s < t < 1\), then \(g_t - g_s = \chi_{[s, t)}\), and \(\lambda(\{x : |\chi_{[s,t)}(x)| > r\}) = t - s > 0\) for every \(r \in (0,1)\), while it equals \(0\) for \(r \ge 1\); hence by the definition of the essential supremum (7.1),

\begin{equation*} \|g_t - g_s\|_\infty = 1 . \end{equation*}

In particular the elements \(g_t\) are distinct elements of \(L^\infty(\mathbb{R})\), so \(S\) is uncountable (the map \(t \mapsto g_t\) from the uncountable set \((0,1)\) into \(S\) is injective), and any two distinct elements of \(S\) are at distance \(1\) from each other.

Now suppose \(C \subseteq L^\infty(\mathbb{R})\) has closure equal to \(L^\infty(\mathbb{R})\). For each \(t \in (0,1)\) choose \(c_t \in C\) with \(\|g_t - c_t\|_\infty < \frac{1}{2}\). If \(c_s = c_t\) with \(s \ne t\), then

\begin{equation*} \|g_s - g_t\|_\infty \le \|g_s - c_s\|_\infty + \|c_t - g_t\|_\infty < 1, \end{equation*}

contradicting \(\|g_s - g_t\|_\infty = 1\). So \(t \mapsto c_t\) is injective from \((0,1)\) into \(C\), forcing \(C\) to be uncountable.

Hence no countable subset of \(L^\infty(\mathbb{R})\) has closure equal to \(L^\infty(\mathbb{R})\).

Problem (7B.4)

Suppose \((X, \mathcal{S}, \mu)\) is a \(\sigma\)-finite measure space and \(1 \le p \le \infty\). Prove that if \(f : X \to \mathbb{F}\) is an \(\mathcal{S}\)-measurable function such that \(fh \in \mathcal{L}^1(\mu)\) for every \(h \in \mathcal{L}^{p^{\prime}}(\mu)\), then \(f \in \mathcal{L}^p(\mu)\).

Solution

Here \(p^{\prime}\) is the dual exponent of \(p\), so \(\frac{1}{p} + \frac{1}{p^{\prime}} = 1\) (see 7.6). The hypothesis says that \(\int |fh| \, d\mu < \infty\) for every \(\mathcal{S}\)-measurable \(h : X \to \mathbb{F}\) with \(\|h\|_{p^{\prime}} < \infty\). Because \(\mu\) is \(\sigma\)-finite, there exist \(X_1 \subseteq X_2 \subseteq \cdots\) in \(\mathcal{S}\) with \(\mu(X_n) < \infty\) for each \(n\) and \(\bigcup_{n=1}^{\infty} X_n = X\) (replace a covering sequence by the increasing sequence of its partial unions).

We first record a small lemma used repeatedly. If \(N \in \mathcal{S}\) with \(\mu(N) = 0\) and \(u : X \to [0, \infty]\) is \(\mathcal{S}\)-measurable, then \(\int u \chi_N \, d\mu = 0\); consequently \(\int u \, d\mu = \int u \chi_{X \setminus N} \, d\mu\). Indeed \(u \chi_N \le \lim_{j \to \infty} j \chi_N\) pointwise, and the Monotone Convergence Theorem (3.11) together with 3.4 gives \(\int \lim_{j} j \chi_N \, d\mu = \lim_j j \mu(N) = 0\); now apply 3.8 and then additivity (3.16) to \(u = u\chi_N + u\chi_{X \setminus N}\).

Case \(p = 1\). Then \(p^{\prime} = \infty\). The constant function \(1\) satisfies \(\|1\|_\infty \le 1 < \infty\), so \(1 \in \mathcal{L}^\infty(\mu)\). By hypothesis \(f \cdot 1 \in \mathcal{L}^1(\mu)\), that is, \(f \in \mathcal{L}^1(\mu)\). (This case does not use \(\sigma\)-finiteness.)

Case \(p = \infty\). Then \(p^{\prime} = 1\). Suppose, for a contradiction, that \(\|f\|_\infty = \infty\). Then for every \(k \in \mathbb{Z}^+\) the set \(A_k = \{x \in X : |f(x)| > 4^k\}\) has \(\mu(A_k) > 0\), by the definition of the essential supremum (7.1). Because \(A_k \cap X_1 \subseteq A_k \cap X_2 \subseteq \cdots\) has union \(A_k\), continuity of measures from below gives \(\mu(A_k \cap X_n) \to \mu(A_k) > 0\), so we may pick \(n_k\) with \(\mu(A_k \cap X_{n_k}) > 0\). Set \(B_k = A_k \cap X_{n_k}\); then \(0 < \mu(B_k) \le \mu(X_{n_k}) < \infty\). Define

\begin{equation*} h = \sum_{k=1}^{\infty} \frac{2^{-k}}{\mu(B_k)} \chi_{B_k}, \end{equation*}

a pointwise sum of nonnegative measurable functions, with values in \([0, \infty]\). By the Monotone Convergence Theorem (3.11) applied to the partial sums,

\begin{equation*} \int h \, d\mu = \sum_{k=1}^{\infty} \frac{2^{-k}}{\mu(B_k)} \mu(B_k) = 1 . \end{equation*}

Hence \(N = \{x : h(x) = \infty\}\) has \(\mu(N) = 0\) (otherwise \(\int h \, d\mu = \infty\) by 3.8, comparing with \(j\chi_N\) and letting \(j \to \infty\)). Let \(\hat{h} = h \chi_{X \setminus N}\), which maps \(X\) into \([0, \infty)\) and satisfies \(\|\hat{h}\|_1 = \int h \, d\mu = 1\) by the lemma; thus \(\hat{h} \in \mathcal{L}^1(\mu) = \mathcal{L}^{p^{\prime}}(\mu)\). But for every \(k\), using the lemma again and \(|f| > 4^k\) on \(B_k\),

\begin{equation*} \int |f \hat{h}| \, d\mu = \int |f| h \, d\mu \ge \frac{2^{-k}}{\mu(B_k)} \int |f| \chi_{B_k} \, d\mu \ge \frac{2^{-k}}{\mu(B_k)} \cdot 4^k \mu(B_k) = 2^k . \end{equation*}

Letting \(k \to \infty\) shows \(\int |f \hat{h}| \, d\mu = \infty\), contradicting the hypothesis that \(f \hat{h} \in \mathcal{L}^1(\mu)\). Hence \(\|f\|_\infty < \infty\), that is, \(f \in \mathcal{L}^\infty(\mu)\).

Case \(1 < p < \infty\). Let

\begin{equation*} E_n = \{x \in X_n : |f(x)| \le n\} \qquad\text{and}\qquad g_n = |f| \chi_{E_n} . \end{equation*}

Then \(E_1 \subseteq E_2 \subseteq \cdots\) and \(\bigcup_n E_n = X\) (every \(x\) lies in some \(X_n\), and \(|f(x)| < \infty\)). Consequently \(g_1^p \le g_2^p \le \cdots\) with \(\lim_{n \to \infty} g_n(x)^p = |f(x)|^p\) for every \(x\), so the Monotone Convergence Theorem (3.11) gives

\begin{equation*} \lim_{n \to \infty} \|g_n\|_p^p = \int |f|^p \, d\mu . \end{equation*}

Note also \(g_n \le n \chi_{X_n}\), so \(\|g_n\|_p^p \le n^p \mu(X_n) < \infty\); thus each \(g_n \in \mathcal{L}^p(\mu)\).

Suppose, for a contradiction, that \(\int |f|^p \, d\mu = \infty\). Then \(\|g_n\|_p \to \infty\), so we may choose \(n_1 < n_2 < \cdots\) with \(c_k := \|g_{n_k}\|_p \ge 4^k\). Write \(G_k = g_{n_k}\), and define

\begin{equation*} h_k = \frac{G_k^{\,p-1}}{c_k^{\,p-1}} . \end{equation*}

Each \(h_k\) is a nonnegative measurable function. Because \((p-1)p^{\prime} = p\) (which follows from \(\frac1p + \frac1{p^{\prime}} = 1\)), we get

\begin{equation*} \|h_k\|_{p^{\prime}}^{p^{\prime}} = \frac{1}{c_k^{\,(p-1)p^{\prime}}} \int G_k^{\,(p-1)p^{\prime}} \, d\mu = \frac{1}{c_k^{\,p}} \int G_k^{\,p} \, d\mu = 1, \end{equation*}

so \(\|h_k\|_{p^{\prime}} = 1\) and \(h_k \in \mathcal{L}^{p^{\prime}}(\mu)\). Furthermore

\begin{equation*} \int G_k h_k \, d\mu = \frac{1}{c_k^{\,p-1}} \int G_k^{\,p} \, d\mu = \frac{c_k^{\,p}}{c_k^{\,p-1}} = c_k \ge 4^k . \end{equation*}

Now define \(h = \sum_{k=1}^{\infty} 2^{-k} h_k\), a pointwise sum of nonnegative measurable functions with values in \([0, \infty]\). Its partial sums \(s_m = \sum_{k=1}^{m} 2^{-k} h_k\) satisfy, by Minkowski’s inequality (7.14),

\begin{equation*} \|s_m\|_{p^{\prime}} \le \sum_{k=1}^{m} 2^{-k} \|h_k\|_{p^{\prime}} \le 1 . \end{equation*}

Since \(s_1^{p^{\prime}} \le s_2^{p^{\prime}} \le \cdots\) converges pointwise to \(h^{p^{\prime}}\), the Monotone Convergence Theorem (3.11) gives \(\int h^{p^{\prime}} \, d\mu = \lim_m \int s_m^{p^{\prime}} \, d\mu \le 1\). In particular the set \(N = \{x : h(x) = \infty\}\) satisfies \(\mu(N) = 0\), so \(\hat{h} = h \chi_{X \setminus N}\) maps \(X\) into \([0, \infty)\) and, by the lemma, \(\|\hat{h}\|_{p^{\prime}} = \|h\|_{p^{\prime}} \le 1\). Thus \(\hat{h} \in \mathcal{L}^{p^{\prime}}(\mu)\).

However, for every \(k\), using \(\hat{h} \ge 2^{-k} h_k \chi_{X \setminus N}\), the lemma, and \(|f| \ge G_k\),

\begin{equation*} \int |f \hat{h}| \, d\mu \ge 2^{-k} \int |f| h_k \chi_{X \setminus N} \, d\mu = 2^{-k} \int |f| h_k \, d\mu \ge 2^{-k} \int G_k h_k \, d\mu \ge 2^{-k} 4^k = 2^k . \end{equation*}

Letting \(k \to \infty\) gives \(\int |f \hat{h}| \, d\mu = \infty\), contradicting the hypothesis that \(f \hat{h} \in \mathcal{L}^1(\mu)\).

Therefore \(\int |f|^p \, d\mu < \infty\), that is, \(f \in \mathcal{L}^p(\mu)\), completing the proof in all cases.

Remark: \(\sigma\)-finiteness cannot be dropped when \(p > 1\). Let \(X = \mathbb{Z}^+\), let \(\mathcal{S}\) be the collection of all subsets of \(X\), and define \(\mu(E) = \infty\) for every nonempty \(E \subseteq X\) and \(\mu(\emptyset) = 0\); this is a measure that is not \(\sigma\)-finite. Because \(p > 1\) we have \(p^{\prime} < \infty\). If \(h : X \to \mathbb{F}\) is not identically \(0\), say \(h(k_0) \ne 0\), then \(\int |h|^{p^{\prime}} \, d\mu \ge |h(k_0)|^{p^{\prime}} \mu(\{k_0\}) = \infty\) by 3.8; thus \(\mathcal{L}^{p^{\prime}}(\mu) = \{0\}\). Hence \(fh = 0 \in \mathcal{L}^1(\mu)\) for every \(\mathcal{S}\)-measurable \(f\), while \(f(k) = k\) satisfies \(\|f\|_p = \infty\). The \(\sigma\)-finiteness hypothesis is used above exactly to produce sets of positive finite measure on which \(|f|\) is large.

Problem (7B.5)

(a) Prove that if \(\mu\) is a measure, \(1 < p < \infty\), and \(f, g \in L^p(\mu)\) are such that

\begin{equation*} \|f\|_p = \|g\|_p = \left\| \frac{f+g}{2} \right\|_p , \end{equation*}

then \(f = g\).

(b) Give an example to show that (a) can fail if \(p = 1\).

(c) Give an example to show that (a) can fail if \(p = \infty\).

Solution

(a) The heart of the matter is a pointwise inequality expressing the strict convexity of \(t \mapsto |t|^p\) for \(1 < p < \infty\).

Lemma. Suppose \(1 < p < \infty\) and \(a, b \in \mathbb{F}\). Then

\begin{equation*} \left| \frac{a+b}{2} \right|^p \le \frac{|a|^p + |b|^p}{2}, \end{equation*}

with equality if and only if \(a = b\).

Proof of Lemma. Let \(\varphi(s) = s^p\) for \(s \in [0, \infty)\). Then \(\varphi\) is strictly increasing, and \(\varphi\) is strictly convex in the following sense: if \(0 \le s < t\) and \(m = \frac{s+t}{2}\), then by the Mean Value Theorem there exist \(\xi_1 \in (s, m)\) and \(\xi_2 \in (m, t)\) with

\begin{equation*} \varphi(m) - \varphi(s) = \varphi^{\prime}(\xi_1)(m - s), \qquad \varphi(t) - \varphi(m) = \varphi^{\prime}(\xi_2)(t - m). \end{equation*}

Because \(\varphi^{\prime}(u) = p u^{p-1}\) is strictly increasing on \([0, \infty)\) (here \(p > 1\) is used) and \(\xi_1 < \xi_2\), and because \(m - s = t - m\), we get \(\varphi(t) - \varphi(m) > \varphi(m) - \varphi(s)\), that is,

\begin{equation*} \varphi\Bigl(\frac{s+t}{2}\Bigr) < \frac{\varphi(s) + \varphi(t)}{2} \qquad\text{whenever } s \ne t . \end{equation*}

Now, using the triangle inequality in \(\mathbb{F}\) and the monotonicity of \(\varphi\),

\begin{equation*} \left| \frac{a+b}{2} \right|^p \le \left( \frac{|a| + |b|}{2} \right)^p \le \frac{|a|^p + |b|^p}{2}, \end{equation*}

where the second inequality is the convexity statement above applied to \(s = |a|\), \(t = |b|\) (with equality when \(|a| = |b|\) and strict inequality when \(|a| \ne |b|\)).

Suppose equality holds overall. Then the second inequality is an equality, so \(|a| = |b| =: r\). Also the first inequality is an equality, and \(\varphi\) is strictly increasing, so \(|a + b| = |a| + |b| = 2r\). If \(r = 0\) then \(a = b = 0\). If \(r > 0\), then

\begin{equation*} 4r^2 = |a+b|^2 = |a|^2 + |b|^2 + 2 \operatorname{Re}(a \overline{b}) = 2r^2 + 2 \operatorname{Re}(a \overline{b}), \end{equation*}

so \(\operatorname{Re}(a \overline{b}) = r^2 = |a| |b| = |a \overline{b}|\). A complex number whose real part equals its absolute value is a nonnegative real number, so \(a \overline{b} = r^2\), which gives

\begin{equation*} a = \frac{r^2}{\overline{b}} = \frac{r^2 b}{|b|^2} = b . \end{equation*}

Conversely if \(a = b\) both sides equal \(|a|^p\). This proves the lemma.

Now suppose \(f, g \in \mathcal{L}^p(\mu)\) represent the given elements of \(L^p(\mu)\) and let

\begin{equation*} c = \|f\|_p = \|g\|_p = \left\| \frac{f+g}{2} \right\|_p < \infty . \end{equation*}

Define \(u : X \to [0, \infty)\) by

\begin{equation*} u = \frac{|f|^p + |g|^p}{2} - \left| \frac{f+g}{2} \right|^p , \end{equation*}

which is nonnegative and measurable by the lemma. By additivity of integration for nonnegative functions (3.16),

\begin{equation*} \int u \, d\mu + \int \left| \frac{f+g}{2} \right|^p d\mu = \int \frac{|f|^p + |g|^p}{2} \, d\mu = \frac{c^p + c^p}{2} = c^p . \end{equation*}

Since \(\int \bigl| \frac{f+g}{2} \bigr|^p d\mu = c^p < \infty\), we may subtract to conclude \(\int u \, d\mu = 0\).

Thus \(u \in \mathcal{L}^1(\mu)\) with \(\|u\|_1 = 0\), so Markov’s inequality (4.1) gives \(\mu(\{x : u(x) \ge \frac{1}{n}\}) \le n \|u\|_1 = 0\) for every \(n \in \mathbb{Z}^+\). Since

\begin{equation*} \{x : u(x) > 0\} = \bigcup_{n=1}^{\infty} \Bigl\{ x : u(x) \ge \tfrac{1}{n} \Bigr\}, \end{equation*}

countable subadditivity gives \(\mu(\{x : u(x) > 0\}) = 0\). Hence \(u(x) = 0\) for almost every \(x \in X\); by the lemma, this means \(f(x) = g(x)\) for almost every \(x \in X\). Therefore \(f\) and \(g\) determine the same element of \(L^p(\mu)\), that is, \(f = g\) in \(L^p(\mu)\).

(b) Let \(\mu\) be counting measure on \(X = \{1, 2\}\), so that \(L^1(\mu) = \mathcal{L}^1(\mu)\) may be identified with \(\mathbb{F}^2\) with the norm \(\|(a_1, a_2)\|_1 = |a_1| + |a_2|\) (counting measure has no nonempty sets of measure \(0\), so no functions are identified). Let

\begin{equation*} f = (1, 0) \qquad\text{and}\qquad g = (0, 1). \end{equation*}

Then \(\|f\|_1 = \|g\|_1 = 1\), while \(\frac{f+g}{2} = (\frac{1}{2}, \frac{1}{2})\) has

\begin{equation*} \left\| \frac{f+g}{2} \right\|_1 = \tfrac{1}{2} + \tfrac{1}{2} = 1 . \end{equation*}

So all three norms equal \(1\), yet \(f \ne g\).

(c) With the same measure space, \(L^\infty(\mu)\) is \(\mathbb{F}^2\) with the norm \(\|(a_1, a_2)\|_\infty = \max\{|a_1|, |a_2|\}\). Let

\begin{equation*} f = (1, 1) \qquad\text{and}\qquad g = (1, 0). \end{equation*}

Then \(\|f\|_\infty = \|g\|_\infty = 1\), while \(\frac{f+g}{2} = (1, \frac{1}{2})\) has \(\bigl\| \frac{f+g}{2} \bigr\|_\infty = 1\). So all three norms equal \(1\), yet \(f \ne g\).

In both (b) and (c) the failure traces back to the lemma: for \(p = 1\) the function \(t \mapsto |t|\) is convex but not strictly convex, and for \(p = \infty\) the norm is not given by an integral of a convex function at all.

Problem (7B.6)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(0 < p < 1\). Show that

\begin{equation*} \|f + g\|_p^p \le \|f\|_p^p + \|g\|_p^p \end{equation*}

for all \(\mathcal{S}\)-measurable functions \(f, g : X \to \mathbb{F}\).

Solution

Both sides may equal \(\infty\); the inequality is to be read in \([0, \infty]\). By the definition of the \(p\)-norm (7.1),

\begin{equation*} \|f + g\|_p^p = \int |f + g|^p \, d\mu, \qquad \|f\|_p^p = \int |f|^p \, d\mu, \qquad \|g\|_p^p = \int |g|^p \, d\mu . \end{equation*}

So it suffices to prove the pointwise inequality \(|f(x) + g(x)|^p \le |f(x)|^p + |g(x)|^p\) for each \(x \in X\) and then integrate.

Step 1: an elementary inequality. If \(a, b \in [0, \infty)\) and \(0 < p < 1\), then

\begin{equation*} (a + b)^p \le a^p + b^p . \end{equation*}

To see this, we may assume \(a + b > 0\) (otherwise both sides are \(0\)). Put

\begin{equation*} s = \frac{a}{a+b}, \qquad t = \frac{b}{a+b}, \end{equation*}

so \(s, t \in [0, 1]\) and \(s + t = 1\). If \(u \in [0, 1]\) and \(0 < p < 1\), then \(u^p \ge u\): this is clear for \(u = 0\), and for \(u \in (0, 1]\) we have \(u^{p} = u \cdot u^{p-1} \ge u\) because \(p - 1 < 0\) and \(u \le 1\) imply \(u^{p-1} \ge 1\). Applying this to \(s\) and to \(t\) gives

\begin{equation*} \frac{a^p + b^p}{(a+b)^p} = s^p + t^p \ge s + t = 1, \end{equation*}

which is the desired inequality.

Step 2: the pointwise estimate. For each \(x \in X\), the triangle inequality in \(\mathbb{F}\) and the fact that \(u \mapsto u^p\) is increasing on \([0, \infty)\) give

\begin{equation*} |f(x) + g(x)|^p \le \bigl( |f(x)| + |g(x)| \bigr)^p \le |f(x)|^p + |g(x)|^p , \end{equation*}

the last step by Step 1 with \(a = |f(x)|\) and \(b = |g(x)|\).

Step 3: integrate. The functions \(|f + g|^p\) and \(|f|^p + |g|^p\) are \(\mathcal{S}\)-measurable functions from \(X\) to \([0, \infty)\), so by the fact that integration is order preserving (3.8) and additivity of integration for nonnegative functions (3.16),

\begin{equation*} \|f + g\|_p^p = \int |f+g|^p \, d\mu \le \int \bigl( |f|^p + |g|^p \bigr) \, d\mu = \int |f|^p \, d\mu + \int |g|^p \, d\mu = \|f\|_p^p + \|g\|_p^p , \end{equation*}

as desired.

Remark: this inequality shows that for \(0 < p < 1\) the formula \(d(f, g) = \|f - g\|_p^p\) satisfies the triangle inequality, so it defines a translation-invariant metric on \(L^p(\mu)\) even though \(\|\cdot\|_p\) itself is not a norm in that range (compare Exercise 7B.1, where the same phenomenon appears on \(\mathbb{F}^n\)).

Problem (7B.7)

Prove that \(L^p(\mu)\), with addition and scalar multiplication as defined in 7.16 and norm defined as in 7.17, is a normed vector space. In other words, prove 7.18.

Solution

The statement to be proved (7.18) is: if \(\mu\) is a measure and \(1 \le p \le \infty\), then \(L^p(\mu)\) is a vector space and \(\|\cdot\|_p\) is a norm on \(L^p(\mu)\).

Throughout, \((X, \mathcal{S}, \mu)\) is a measure space and \(1 \le p \le \infty\). Recall from 7.15 that \(Z(\mu)\) is the set of \(\mathcal{S}\)-measurable functions from \(X\) to \(\mathbb{F}\) that equal \(0\) almost everywhere, that \(\widetilde{f} = \{f + z : z \in Z(\mu)\}\) for \(f \in \mathcal{L}^p(\mu)\), and from 7.16 that \(L^p(\mu) = \{\widetilde{f} : f \in \mathcal{L}^p(\mu)\}\) with

\begin{equation*} \widetilde{f} + \widetilde{g} = (f+g)^{\sim} \qquad\text{and}\qquad \alpha \widetilde{f} = (\alpha f)^{\sim} . \end{equation*}

We use the following lemma twice: if \(N \in \mathcal{S}\) with \(\mu(N) = 0\) and \(u : X \to [0, \infty]\) is \(\mathcal{S}\)-measurable, then \(\int u \chi_N \, d\mu = 0\), and hence \(\int u \, d\mu = \int u \chi_{X \setminus N} \, d\mu\). Indeed \(u \chi_N \le \lim_{j \to \infty} j \chi_N\) pointwise, and by the Monotone Convergence Theorem (3.11) together with 3.4 the right side has integral \(\lim_j j \mu(N) = 0\); now use 3.8 and additivity (3.16).

Step 1: \(\mathcal{L}^p(\mu)\) is a vector space and \(Z(\mu)\) is a subspace of it. For \(1 \le p < \infty\) this is 7.5. For \(p = \infty\), the inequalities \(\|f+g\|_\infty \le \|f\|_\infty + \|g\|_\infty\) and \(\|\alpha f\|_\infty = |\alpha| \|f\|_\infty\) (Exercise 1 in Section 7A) show that \(\mathcal{L}^\infty(\mu)\) is closed under addition and scalar multiplication, hence is a subspace of the vector space of all \(\mathcal{S}\)-measurable functions from \(X\) to \(\mathbb{F}\) with pointwise operations. As noted after 7.15, \(Z(\mu)\) is closed under scalar multiplication, and it is closed under addition because the union of two sets of measure \(0\) has measure \(0\); also \(Z(\mu) \subseteq \mathcal{L}^p(\mu)\) because a function that vanishes almost everywhere has \(\|\cdot\|_p = 0\). Thus \(Z(\mu)\) is a subspace of \(\mathcal{L}^p(\mu)\).

Step 2: \(\widetilde{f} = \widetilde{F}\) if and only if \(f - F \in Z(\mu)\), if and only if \(f = F\) almost everywhere. If \(f - F \in Z(\mu)\), then for \(z \in Z(\mu)\) we have \(f + z = F + \bigl( (f - F) + z \bigr) \in \widetilde{F}\) and similarly \(F + z \in \widetilde{f}\), because \(Z(\mu)\) is closed under addition and under multiplication by \(-1\); hence \(\widetilde{f} = \widetilde{F}\). Conversely if \(\widetilde{f} = \widetilde{F}\), then \(f = f + 0 \in \widetilde{f} = \widetilde{F}\), so \(f = F + z\) for some \(z \in Z(\mu)\), giving \(f - F \in Z(\mu)\). The second equivalence is the definition of \(Z(\mu)\).

Step 3: the operations in 7.16 are well defined. Suppose \(f, F, g, G \in \mathcal{L}^p(\mu)\) with \(\widetilde{f} = \widetilde{F}\) and \(\widetilde{g} = \widetilde{G}\), and suppose \(\alpha \in \mathbb{F}\). By Step 2, \(f - F \in Z(\mu)\) and \(g - G \in Z(\mu)\). Since \(Z(\mu)\) is a subspace,

\begin{equation*} (f+g) - (F+G) = (f - F) + (g - G) \in Z(\mu) \qquad\text{and}\qquad \alpha f - \alpha F = \alpha (f - F) \in Z(\mu), \end{equation*}

so by Step 2 again \((f+g)^{\sim} = (F+G)^{\sim}\) and \((\alpha f)^{\sim} = (\alpha F)^{\sim}\). Thus the sum \(\widetilde{f} + \widetilde{g}\) and the scalar multiple \(\alpha \widetilde{f}\) do not depend on the representatives chosen.

Step 4: \(L^p(\mu)\) is a vector space. Every element of \(L^p(\mu)\) is of the form \(\widetilde{f}\) with \(f \in \mathcal{L}^p(\mu)\), and by definition the map \(f \mapsto \widetilde{f}\) carries sums to sums and scalar multiples to scalar multiples. Hence each vector space axiom for \(L^p(\mu)\) follows from the corresponding axiom in \(\mathcal{L}^p(\mu)\). Explicitly, for \(f, g, k \in \mathcal{L}^p(\mu)\) and \(\alpha, \beta \in \mathbb{F}\):

\begin{equation*} \widetilde{f} + \widetilde{g} = (f+g)^{\sim} = (g+f)^{\sim} = \widetilde{g} + \widetilde{f}, \end{equation*}

\begin{equation*} (\widetilde{f} + \widetilde{g}) + \widetilde{k} = \bigl( (f+g) + k \bigr)^{\sim} = \bigl( f + (g+k) \bigr)^{\sim} = \widetilde{f} + (\widetilde{g} + \widetilde{k}), \end{equation*}

\begin{equation*} \widetilde{f} + \widetilde{0} = (f + 0)^{\sim} = \widetilde{f}, \qquad \widetilde{f} + (-f)^{\sim} = (f - f)^{\sim} = \widetilde{0}, \end{equation*}

\begin{equation*} 1 \widetilde{f} = (1f)^{\sim} = \widetilde{f}, \qquad \alpha(\beta \widetilde{f}) = (\alpha \beta f)^{\sim} = (\alpha\beta) \widetilde{f}, \end{equation*}

\begin{equation*} \alpha(\widetilde{f} + \widetilde{g}) = \bigl( \alpha(f+g) \bigr)^{\sim} = (\alpha f + \alpha g)^{\sim} = \alpha \widetilde{f} + \alpha \widetilde{g}, \qquad (\alpha + \beta) \widetilde{f} = (\alpha f + \beta f)^{\sim} = \alpha \widetilde{f} + \beta \widetilde{f} . \end{equation*}

Thus \(L^p(\mu)\) is a vector space over \(\mathbb{F}\), with additive identity \(\widetilde{0} = Z(\mu)\) and with \(-\widetilde{f} = (-f)^{\sim}\).

Step 5: \(\|\cdot\|_p\) is well defined on \(L^p(\mu)\). Suppose \(f, F \in \mathcal{L}^p(\mu)\) with \(\widetilde{f} = \widetilde{F}\); by Step 2 there is \(N \in \mathcal{S}\) with \(\mu(N) = 0\) and \(f = F\) on \(X \setminus N\). If \(p < \infty\), then \(|f|^p \chi_{X \setminus N} = |F|^p \chi_{X \setminus N}\), so the lemma gives

\begin{equation*} \int |f|^p \, d\mu = \int |f|^p \chi_{X \setminus N} \, d\mu = \int |F|^p \chi_{X \setminus N} \, d\mu = \int |F|^p \, d\mu, \end{equation*}

hence \(\|f\|_p = \|F\|_p\). If \(p = \infty\), then for each \(t > 0\) the sets \(\{x : |f(x)| > t\}\) and \(\{x : |F(x)| > t\}\) differ only inside \(N\), so one has measure \(0\) if and only if the other does; by the definition of the essential supremum (7.1) the two infima agree, so \(\|f\|_\infty = \|F\|_\infty\). Thus \(\|\widetilde{f}\|_p = \|f\|_p\) is unambiguous. Moreover \(\|\widetilde{f}\|_p < \infty\) because \(f \in \mathcal{L}^p(\mu)\), so \(\|\cdot\|_p\) maps \(L^p(\mu)\) into \([0, \infty)\), as required by 6.33.

Step 6: the three norm properties of 6.33.

Positive definiteness. Suppose \(\|\widetilde{f}\|_p = 0\). If \(p < \infty\), then \(|f|^p\) is a nonnegative function with \(\int |f|^p \, d\mu = 0\), so \(|f|^p \in \mathcal{L}^1(\mu)\) and Markov’s inequality (4.1) gives

\begin{equation*} \mu\Bigl( \bigl\{ x : |f(x)|^p \ge \tfrac{1}{n} \bigr\} \Bigr) \le n \int |f|^p \, d\mu = 0 \end{equation*}

for every \(n \in \mathbb{Z}^+\); since \(\{x : f(x) \ne 0\}\) is the union over \(n \in \mathbb{Z}^+\) of these sets, countable subadditivity gives \(\mu(\{x : f(x) \ne 0\}) = 0\). If \(p = \infty\), then \(\|f\|_\infty = 0\) means that for each \(n \in \mathbb{Z}^+\) there is \(t \in (0, \frac{1}{n})\) with \(\mu(\{x : |f(x)| > t\}) = 0\), and hence \(\mu(\{x : |f(x)| > \frac{1}{n}\}) = 0\); again taking the union over \(n\) gives \(\mu(\{x : f(x) \ne 0\}) = 0\). In both cases \(f \in Z(\mu)\), so \(\widetilde{f} = \widetilde{0}\) by Step 2. Conversely, if \(\widetilde{f} = \widetilde{0}\) then \(f = 0\) almost everywhere, and Step 5 (applied with \(F = 0\)) gives \(\|\widetilde{f}\|_p = \|0\|_p = 0\).

Homogeneity. For \(\alpha \in \mathbb{F}\) and \(f \in \mathcal{L}^p(\mu)\),

\begin{equation*} \|\alpha \widetilde{f}\|_p = \|(\alpha f)^{\sim}\|_p = \|\alpha f\|_p = |\alpha| \, \|f\|_p = |\alpha| \, \|\widetilde{f}\|_p, \end{equation*}

where the middle equality is 7.5(b) if \(p < \infty\) and is Exercise 1 in Section 7A if \(p = \infty\).

Triangle inequality. For \(f, g \in \mathcal{L}^p(\mu)\), Minkowski’s inequality (7.14) gives

\begin{equation*} \|\widetilde{f} + \widetilde{g}\|_p = \|(f+g)^{\sim}\|_p = \|f + g\|_p \le \|f\|_p + \|g\|_p = \|\widetilde{f}\|_p + \|\widetilde{g}\|_p . \end{equation*}

Thus \(\|\cdot\|_p\) satisfies all three conditions in the definition of a norm (6.33), so \(L^p(\mu)\) with the operations of 7.16 and the norm of 7.17 is a normed vector space. This proves 7.18.

Problem (7B.8)

Prove 7.20 for the case \(p = \infty\).

Solution

We must prove the following statement. Suppose \((X, \mathcal{S}, \mu)\) is a measure space and \(f_1, f_2, \ldots\) is a sequence of functions in \(\mathcal{L}^\infty(\mu)\) such that for every \(\varepsilon > 0\) there exists \(n \in \mathbf{Z}^+\) with \(\|f_j - f_k\|_\infty < \varepsilon\) for all \(j \ge n\) and \(k \ge n\). Then there exists \(f \in \mathcal{L}^\infty(\mu)\) such that \(\lim_{k \to \infty} \|f_k - f\|_\infty = 0\).

We begin with a small lemma that we use repeatedly.

Lemma. If \(g : X \to \mathbf{F}\) is \(\mathcal{S}\)-measurable, then

\begin{equation*} \mu\bigl(\{x \in X : |g(x)| > \|g\|_\infty\}\bigr) = 0 . \end{equation*}

Proof of Lemma. We may assume \(\|g\|_\infty < \infty\). For each \(m \in \mathbf{Z}^+\), the definition of the essential supremum (7.1) as an infimum produces \(t_m > 0\) with \(t_m < \|g\|_\infty + \frac{1}{m}\) and \(\mu(\{x : |g(x)| > t_m\}) = 0\). Since

\begin{equation*} \Bigl\{x : |g(x)| > \|g\|_\infty + \tfrac{1}{m}\Bigr\} \subseteq \{x : |g(x)| > t_m\}, \end{equation*}

each of the sets on the left has \(\mu\)-measure \(0\). Because

\begin{equation*} \{x : |g(x)| > \|g\|_\infty\} = \bigcup_{m=1}^{\infty} \Bigl\{x : |g(x)| > \|g\|_\infty + \tfrac{1}{m}\Bigr\}, \end{equation*}

countable subadditivity of \(\mu\) (2.58) gives the conclusion. Note also that each of these sets lies in \(\mathcal{S}\), because \(|g|\) is \(\mathcal{S}\)-measurable.

Now we prove the result. For \(j, k \in \mathbf{Z}^+\) set

\begin{equation*} E_{j,k} = \{x \in X : |f_j(x) - f_k(x)| > \|f_j - f_k\|_\infty\} . \end{equation*}

Each \(E_{j,k}\) is in \(\mathcal{S}\) and, by the Lemma applied to \(g = f_j - f_k\), satisfies \(\mu(E_{j,k}) = 0\). Also let \(F_k = \{x \in X : |f_k(x)| > \|f_k\|_\infty\}\), which similarly lies in \(\mathcal{S}\) and has measure \(0\). Put

\begin{equation*} E = \Bigl(\bigcup_{j,k \in \mathbf{Z}^+} E_{j,k}\Bigr) \cup \Bigl(\bigcup_{k \in \mathbf{Z}^+} F_k\Bigr). \end{equation*}

This is a countable union of sets of measure \(0\), so \(E \in \mathcal{S}\) and \(\mu(E) = 0\).

If \(x \in X \setminus E\), then

\begin{equation*} |f_j(x) - f_k(x)| \le \|f_j - f_k\|_\infty \quad \text{for all } j, k \in \mathbf{Z}^+ . \end{equation*}

Hence, given \(\varepsilon > 0\), choosing \(n \in \mathbf{Z}^+\) as in the hypothesis, we get

\begin{equation*} |f_j(x) - f_k(x)| < \varepsilon \quad \text{for all } j \ge n,\ k \ge n \text{ and all } x \in X \setminus E . \end{equation*}

In other words, \(f_1, f_2, \ldots\) is uniformly Cauchy on \(X \setminus E\). In particular, for each \(x \in X \setminus E\) the sequence \(f_1(x), f_2(x), \ldots\) is a Cauchy sequence in \(\mathbf{F}\), and therefore converges, because \(\mathbf{R}\) and \(\mathbf{C}\) are complete.

Define \(f : X \to \mathbf{F}\) by

\begin{equation*} f(x) = \begin{cases} \lim_{k \to \infty} f_k(x) & \text{if } x \in X \setminus E, \\ 0 & \text{if } x \in E. \end{cases} \end{equation*}

Each function \(f_k \chi_{X \setminus E}\) is \(\mathcal{S}\)-measurable (a product of \(\mathcal{S}\)-measurable functions), and \(f(x) = \lim_{k \to \infty} \bigl(f_k \chi_{X \setminus E}\bigr)(x)\) for every \(x \in X\). Thus \(f\) is \(\mathcal{S}\)-measurable by 2.48 (applied to the real and imaginary parts if \(\mathbf{F} = \mathbf{C}\)).

Now fix \(\varepsilon > 0\) and choose \(n \in \mathbf{Z}^+\) such that \(\|f_j - f_k\|_\infty < \varepsilon\) for all \(j, k \ge n\). Suppose \(k \ge n\) and \(x \in X \setminus E\). Then

\begin{equation*} |f_k(x) - f(x)| = \lim_{j \to \infty} |f_k(x) - f_j(x)| \le \varepsilon . \end{equation*}

Therefore

\begin{equation*} \{x \in X : |f_k(x) - f(x)| > \varepsilon\} \subseteq E , \end{equation*}

which has measure \(0\). Consequently \(\mu(\{x : |f_k(x) - f(x)| > t\}) = 0\) for every \(t > \varepsilon\), and hence \(\|f_k - f\|_\infty \le \varepsilon\) for all \(k \ge n\).

It remains to check that \(f \in \mathcal{L}^\infty(\mu)\). Take \(\varepsilon = 1\) and let \(n\) be as above. For \(x \in X \setminus E\) we have \(|f(x)| \le |f_n(x)| + |f(x) - f_n(x)| \le \|f_n\|_\infty + 1\) (the first bound because \(x \notin F_n\)). Thus \(\{x : |f(x)| > \|f_n\|_\infty + 1\} \subseteq E\) has measure \(0\), so

\begin{equation*} \|f\|_\infty \le \|f_n\|_\infty + 1 < \infty , \end{equation*}

which shows \(f \in \mathcal{L}^\infty(\mu)\).

Finally, the estimate \(\|f_k - f\|_\infty \le \varepsilon\) for all \(k \ge n\), valid for each \(\varepsilon > 0\) with a suitable \(n\), says precisely that \(\lim_{k \to \infty} \|f_k - f\|_\infty = 0\), completing the proof.

Problem (7B.9)

Prove that 7.20 also holds for \(p \in (0, 1)\).

Solution

We must prove: if \((X, \mathcal{S}, \mu)\) is a measure space, \(0 < p < 1\), and \(f_1, f_2, \ldots\) is a sequence in \(\mathcal{L}^p(\mu)\) such that for every \(\varepsilon > 0\) there exists \(n \in \mathbf{Z}^+\) with \(\|f_j - f_k\|_p < \varepsilon\) for all \(j \ge n\) and \(k \ge n\), then there exists \(f \in \mathcal{L}^p(\mu)\) with \(\lim_{k \to \infty}\|f_k - f\|_p = 0\).

Throughout we use Exercise 6 in this section, which states that for \(0 < p < 1\) and all \(\mathcal{S}\)-measurable \(u, v : X \to \mathbf{F}\),

\begin{equation*} \|u + v\|_p^p \le \|u\|_p^p + \|v\|_p^p . \end{equation*}

An easy induction extends this to finitely many summands: \(\bigl\|\sum_{k=1}^m u_k\bigr\|_p^p \le \sum_{k=1}^m \|u_k\|_p^p\). Note that \(\|u\|_p^p = \int |u|^p \, d\mu\), and that because \(t \mapsto t^p\) is a strictly increasing bijection of \([0, \infty)\) onto itself, the hypothesis is equivalent to: for every \(\varepsilon > 0\) there is \(n\) with \(\|f_j - f_k\|_p^p < \varepsilon\) for all \(j, k \ge n\).

Step 1: it suffices to find a subsequence that converges. Suppose \(f \in \mathcal{L}^p(\mu)\) and \(f_{k_1}, f_{k_2}, \ldots\) is a subsequence with \(\lim_{m \to \infty} \|f_{k_m} - f\|_p = 0\). Let \(\varepsilon > 0\). Choose \(n\) with \(\|f_j - f_k\|_p^p < \varepsilon/2\) for all \(j, k \ge n\), and then choose \(m\) with \(k_m \ge n\) and \(\|f_{k_m} - f\|_p^p < \varepsilon/2\). For every \(k \ge n\), Exercise 6 gives

\begin{equation*} \|f_k - f\|_p^p \le \|f_k - f_{k_m}\|_p^p + \|f_{k_m} - f\|_p^p < \varepsilon . \end{equation*}

Hence \(\lim_{k \to \infty}\|f_k - f\|_p = 0\). (This is the analogue of Exercise 14 in Section 6A; only the triangle inequality for \(\|\cdot\|_p^p\) was used.)

Step 2: choice of subsequence. For each \(m \in \mathbf{Z}^+\) use the hypothesis to pick \(n_m \in \mathbf{Z}^+\) with \(\|f_j - f_k\|_p^p < 2^{-m}\) for all \(j, k \ge n_m\); replacing \(n_m\) by \(\max\{n_1, \ldots, n_m\} + m\) we may assume \(n_1 < n_2 < \cdots\). Setting \(k_m = n_m\) we obtain a subsequence with

\begin{equation*} \sum_{m=2}^{\infty} \|f_{k_m} - f_{k_{m-1}}\|_p^p \le \sum_{m=2}^{\infty} 2^{-(m-1)} < \infty . \end{equation*}

A subsequence of our sequence still satisfies the Cauchy hypothesis. Thus, dropping to this subsequence but not relabeling, and setting \(f_0 = 0\), we may assume

\begin{equation*} \sum_{k=1}^{\infty} \|f_k - f_{k-1}\|_p^p < \infty . \end{equation*}

Step 3: pointwise limit. Define \(g_m, g : X \to [0, \infty]\) by

\begin{equation*} g_m(x) = \sum_{k=1}^{m} |f_k(x) - f_{k-1}(x)| , \qquad g(x) = \sum_{k=1}^{\infty} |f_k(x) - f_{k-1}(x)| . \end{equation*}

By the finite form of Exercise 6,

\begin{equation*} \int g_m^{\,p} \, d\mu = \|g_m\|_p^p \le \sum_{k=1}^{m} \|f_k - f_{k-1}\|_p^p \le \sum_{k=1}^{\infty} \|f_k - f_{k-1}\|_p^p . \end{equation*}

Since \(0 \le g_1 \le g_2 \le \cdots\) and \(g_m \to g\) pointwise, also \(g_m^{\,p} \uparrow g^p\) pointwise (because \(t \mapsto t^p\) is increasing and continuous on \([0, \infty]\)). The Monotone Convergence Theorem (3.11) therefore gives

\begin{equation*} \int g^p \, d\mu = \lim_{m \to \infty} \int g_m^{\,p} \, d\mu \le \sum_{k=1}^{\infty} \|f_k - f_{k-1}\|_p^p < \infty . \end{equation*}

Hence the \(\mathcal{S}\)-measurable set \(A = \{x \in X : g(x) < \infty\}\) satisfies \(\mu(X \setminus A) = 0\). Indeed, \(X \setminus A = \{x \in X : g(x) = \infty\}\), so for every \(c > 0\) we have \(c \chi_{X \setminus A} \le g^p\) pointwise on \(X\), whence \(c \, \mu(X \setminus A) \le \int g^p \, d\mu < \infty\). Letting \(c \to \infty\) forces \(\mu(X \setminus A) = 0\).

For \(x \in A\) the series \(\sum_{k=1}^\infty \bigl(f_k(x) - f_{k-1}(x)\bigr)\) converges absolutely, hence converges, and its \(m\)-th partial sum is \(f_m(x)\). Thus \(\lim_{m \to \infty} f_m(x)\) exists for every \(x \in A\). Define

\begin{equation*} f(x) = \lim_{m \to \infty} \bigl(f_m \chi_A\bigr)(x) , \end{equation*}

so that \(f = \lim_m f_m\) on \(A\) and \(f = 0\) on \(X \setminus A\). Each \(f_m \chi_A\) is \(\mathcal{S}\)-measurable and the limit exists at every point of \(X\), so \(f\) is \(\mathcal{S}\)-measurable by 2.48 (applied to the real and imaginary parts if \(\mathbf{F} = \mathbf{C}\)). Moreover \(|f| \le g\) on \(A\), hence \(|f|^p \le g^p\) except on the null set \(X \setminus A\), and therefore

\begin{equation*} \|f\|_p^p = \int |f|^p \, d\mu \le \int g^p \, d\mu < \infty , \end{equation*}

so \(f \in \mathcal{L}^p(\mu)\).

Step 4: convergence in \(\|\cdot\|_p\). Let \(\varepsilon > 0\) and choose \(n\) with \(\|f_j - f_k\|_p^p < \varepsilon\) for all \(j, k \ge n\). Fix \(k \ge n\). For every \(x \in A\) we have \(|f_k(x) - f_j(x)|^p \to |f_k(x) - f(x)|^p\) as \(j \to \infty\), so this convergence holds almost everywhere on \(X\). Fatou’s Lemma (Exercise 17 in Section 3A) applied to the nonnegative measurable functions \(|f_k - f_j|^p\) gives

\begin{equation*} \|f_k - f\|_p^p = \int |f_k - f|^p \, d\mu \le \liminf_{j \to \infty} \int |f_k - f_j|^p \, d\mu = \liminf_{j \to \infty} \|f_k - f_j\|_p^p \le \varepsilon . \end{equation*}

Hence \(\lim_{k \to \infty} \|f_k - f\|_p = 0\) for the subsequence to which we passed in Step 2, and Step 1 upgrades this to the original sequence. This proves 7.20 for \(p \in (0,1)\).

Problem (7B.10)

Prove that 7.23 also holds for \(p \in (0, 1)\).

Solution

We must prove: if \((X, \mathcal{S}, \mu)\) is a measure space, \(0 < p < 1\), \(f \in \mathcal{L}^p(\mu)\), and \(f_1, f_2, \ldots\) is a sequence in \(\mathcal{L}^p(\mu)\) with \(\lim_{k \to \infty} \|f_k - f\|_p = 0\), then there is a subsequence \(f_{k_1}, f_{k_2}, \ldots\) such that \(\lim_{m \to \infty} f_{k_m}(x) = f(x)\) for almost every \(x \in X\).

Because \(\|f_k - f\|_p \to 0\), we also have \(\|f_k - f\|_p^p = \int |f_k - f|^p \, d\mu \to 0\). Hence we may choose \(k_1 < k_2 < \cdots\) in \(\mathbf{Z}^+\) such that

\begin{equation*} \int |f_{k_m} - f|^p \, d\mu < 4^{-m} \qquad \text{for every } m \in \mathbf{Z}^+ . \end{equation*}

For each \(m \in \mathbf{Z}^+\) put \(h_m = |f_{k_m} - f|^p\). Each \(h_m\) is a nonnegative \(\mathcal{S}\)-measurable function with \(\|h_m\|_1 = \int h_m \, d\mu < 4^{-m} < \infty\), so \(h_m \in \mathcal{L}^1(\mu)\). Define

\begin{equation*} E_m = \{x \in X : h_m(x) \ge 2^{-m}\} = \{x \in X : |f_{k_m}(x) - f(x)|^p \ge 2^{-m}\} . \end{equation*}

These sets belong to \(\mathcal{S}\), and Markov’s inequality (4.1), applied with \(c = 2^{-m}\), gives

\begin{equation*} \mu(E_m) \le 2^{m} \|h_m\|_1 < 2^m \cdot 4^{-m} = 2^{-m} . \end{equation*}

For \(M \in \mathbf{Z}^+\) let \(F_M = \bigcup_{m=M}^{\infty} E_m\), and let \(F = \bigcap_{M=1}^{\infty} F_M\). Since \(F \subseteq F_M\), monotonicity of \(\mu\) (2.57) and countable subadditivity (2.58) give

\begin{equation*} \mu(F) \le \mu(F_M) \le \sum_{m=M}^{\infty} \mu(E_m) < \sum_{m=M}^{\infty} 2^{-m} = 2^{-M+1} \end{equation*}

for every \(M \in \mathbf{Z}^+\). Letting \(M \to \infty\) shows \(\mu(F) = 0\).

Now suppose \(x \in X \setminus F\). Then \(x \notin F_M\) for some \(M \in \mathbf{Z}^+\), which means \(x \notin E_m\) for every \(m \ge M\); that is,

\begin{equation*} |f_{k_m}(x) - f(x)|^p < 2^{-m} \qquad \text{for all } m \ge M . \end{equation*}

Therefore \(|f_{k_m}(x) - f(x)| < 2^{-m/p} \to 0\) as \(m \to \infty\), so \(\lim_{m \to \infty} f_{k_m}(x) = f(x)\).

Thus \(\lim_{m \to \infty} f_{k_m}(x) = f(x)\) for every \(x \in X \setminus F\), and \(\mu(F) = 0\). Hence the subsequence \(f_{k_1}, f_{k_2}, \ldots\) converges to \(f\) pointwise almost everywhere, as required.

Note that nothing in this argument used \(p < 1\): the same proof gives 7.23 for every \(p \in (0, \infty)\), since the only tools were Markov’s inequality and countable subadditivity. (For \(p = \infty\) the conclusion of 7.23 holds for the full sequence, since \(\|f_k - f\|_\infty \to 0\) forces \(f_k \to f\) uniformly off a set of measure \(0\), by the Lemma proved in Exercise 8.)

Problem (7B.11)

Suppose \(1 \le p \le \infty\). Prove that

\begin{equation*} \{(a_1, a_2, \ldots) \in \ell^p : a_k \ne 0 \text{ for every } k \in \mathbf{Z}^+\} \end{equation*}

is not an open subset of \(\ell^p\).

Solution

Let

\begin{equation*} A = \{(a_1, a_2, \ldots) \in \ell^p : a_k \ne 0 \text{ for every } k \in \mathbf{Z}^+\} . \end{equation*}

Recall that a subset of a normed vector space is open exactly when each of its elements is an interior point of it. So it suffices to exhibit one element \(a \in A\) such that no open ball centered at \(a\) is contained in \(A\).

Let \(a = (a_1, a_2, \ldots)\) where \(a_k = 2^{-k}\). Then \(a_k \ne 0\) for every \(k\), and \(a \in \ell^p\) for every \(p\) with \(1 \le p \le \infty\): indeed, if \(p < \infty\) then

\begin{equation*} \|a\|_p^p = \sum_{k=1}^{\infty} 2^{-kp} = \frac{2^{-p}}{1 - 2^{-p}} < \infty , \end{equation*}

and \(\|a\|_\infty = \sup\{2^{-k} : k \in \mathbf{Z}^+\} = \tfrac12 < \infty\). Thus \(a \in A\).

Now let \(\varepsilon > 0\). Choose \(n \in \mathbf{Z}^+\) with \(2^{-n} < \varepsilon\). Define \(b = (b_1, b_2, \ldots)\) by

\begin{equation*} b_k = \begin{cases} 2^{-k} & \text{if } k \ne n, \\ 0 & \text{if } k = n. \end{cases} \end{equation*}

Then \(|b_k| \le |a_k|\) for every \(k\), so \(b \in \ell^p\). Furthermore \(a - b = 2^{-n} e_n\), where \(e_n\) is the sequence whose \(n\)-th term is \(1\) and whose other terms are \(0\). Hence, for every \(p\) with \(1 \le p \le \infty\),

\begin{equation*} \|a - b\|_p = 2^{-n}\|e_n\|_p = 2^{-n} < \varepsilon \end{equation*}

(this holds for \(p < \infty\) because \(\|e_n\|_p = (1^p)^{1/p} = 1\), and for \(p = \infty\) because \(\|e_n\|_\infty = 1\); recall from 7.2 that for counting measure on \(\mathbf{Z}^+\) the essential supremum is the supremum).

Since \(b_n = 0\), we have \(b \notin A\). So \(b\) belongs to the open ball of radius \(\varepsilon\) centered at \(a\) but not to \(A\). As \(\varepsilon > 0\) was arbitrary, no open ball centered at \(a\) is contained in \(A\), and therefore \(a\) is not an interior point of \(A\).

Because \(a \in A\) but \(a\) is not an interior point of \(A\), the set \(A\) is not open in \(\ell^p\).

Remark. For \(1 \le p < \infty\) this argument works starting from any \(a \in A\): since \(\sum_k |a_k|^p < \infty\) we have \(a_k \to 0\), so arbitrarily small perturbations that zero out a single far-out coordinate leave \(A\). Thus for \(p < \infty\) the set \(A\) has empty interior. This fails for \(p = \infty\), where for instance \((1, 1, 1, \ldots)\) is an interior point of \(A\) (every \(b\) with \(\|b - (1,1,1,\ldots)\|_\infty < \tfrac12\) satisfies \(|b_k| > \tfrac12\) for all \(k\)). This is why the specific choice of \(a\) above matters when \(p = \infty\).

Problem (7B.12)

Show that there exists a sequence \(f_1, f_2, \ldots\) of functions in \(\mathcal{L}^1([0,1])\) such that \(\lim_{k \to \infty} \|f_k\|_1 = 0\) but

\begin{equation*} \sup\{f_k(x) : k \in \mathbf{Z}^+\} = \infty \end{equation*}

for every \(x \in [0,1]\).

[This exercise shows that the conclusion of 7.23 cannot be improved to conclude that \(\lim_{k \to \infty} f_k(x) = f(x)\) for almost every \(x \in X\).]

Solution

For \(n \in \mathbf{Z}^+\) and \(j \in \{0, 1, \ldots, 2^n - 1\}\), let

\begin{equation*} I_{n,j} = \Bigl[\frac{j}{2^n}, \frac{j+1}{2^n}\Bigr] \subseteq [0,1] . \end{equation*}

For each fixed \(n\), these \(2^n\) closed intervals cover \([0,1]\).

We index all such pairs \((n,j)\) by a single positive integer. The number of pairs with first coordinate less than \(n\) is \(\sum_{i=1}^{n-1} 2^i = 2^n - 2\). Hence every \(k \in \mathbf{Z}^+\) can be written uniquely as

\begin{equation*} k = 2^n - 1 + j \qquad \text{with } n \in \mathbf{Z}^+ \text{ and } j \in \{0, 1, \ldots, 2^n - 1\} , \end{equation*}

because as \(j\) runs over \(\{0, \ldots, 2^n-1\}\) the integer \(k\) runs over \(\{2^n - 1, \ldots, 2^{n+1} - 2\}\), and these blocks partition \(\mathbf{Z}^+\) as \(n\) runs over \(\mathbf{Z}^+\). Write \(n(k)\) and \(j(k)\) for the resulting \(n\) and \(j\). Now define

\begin{equation*} f_k = n(k) \, \chi_{I_{n(k),\, j(k)}} . \end{equation*}

Thus the sequence begins

\begin{equation*} 1 \cdot \chi_{[0,\frac12]},\ 1 \cdot \chi_{[\frac12,1]},\ 2 \cdot \chi_{[0,\frac14]},\ 2 \cdot \chi_{[\frac14,\frac12]},\ 2 \cdot \chi_{[\frac12,\frac34]},\ 2 \cdot \chi_{[\frac34,1]},\ 3 \cdot \chi_{[0,\frac18]},\ \ldots \end{equation*}

Each \(f_k\) is a nonnegative Borel measurable simple function on \([0,1]\), and

\begin{equation*} \|f_k\|_1 = \int_{[0,1]} f_k \, d\lambda = n(k) \, \lambda\bigl(I_{n(k), j(k)}\bigr) = \frac{n(k)}{2^{n(k)}} < \infty , \end{equation*}

so \(f_k \in \mathcal{L}^1([0,1])\) for every \(k\).

We check that \(\|f_k\|_1 \to 0\). From \(k = 2^{n(k)} - 1 + j(k) \ge 2^{n(k)} - 1\) we get \(n(k) \le \log_2(k+1)\), and from \(k \le 2^{n(k)+1} - 2\) we get \(n(k) \ge \log_2(k+2) - 1\). In particular \(n(k) \to \infty\) as \(k \to \infty\). Since \(\lim_{n \to \infty} n 2^{-n} = 0\), given \(\varepsilon > 0\) there is \(N\) with \(n 2^{-n} < \varepsilon\) for all \(n \ge N\); choosing \(K\) so that \(n(k) \ge N\) whenever \(k \ge K\), we get \(\|f_k\|_1 = n(k) 2^{-n(k)} < \varepsilon\) for all \(k \ge K\). Hence

\begin{equation*} \lim_{k \to \infty} \|f_k\|_1 = 0 . \end{equation*}

Finally, fix \(x \in [0,1]\) and let \(n \in \mathbf{Z}^+\). Since \(I_{n,0}, \ldots, I_{n, 2^n - 1}\) cover \([0,1]\), there is \(j \in \{0, \ldots, 2^n - 1\}\) with \(x \in I_{n,j}\). For \(k = 2^n - 1 + j\) we then have \(n(k) = n\), \(j(k) = j\), and therefore

\begin{equation*} f_k(x) = n \, \chi_{I_{n,j}}(x) = n . \end{equation*}

Thus \(\sup\{f_k(x) : k \in \mathbf{Z}^+\} \ge n\) for every \(n \in \mathbf{Z}^+\), which gives

\begin{equation*} \sup\{f_k(x) : k \in \mathbf{Z}^+\} = \infty \end{equation*}

for every \(x \in [0,1]\), as desired.

Remark. Here \(\|f_k - 0\|_1 \to 0\), so \(f_k \to 0\) in \(L^1([0,1])\), and yet \(\lim_{k \to \infty} f_k(x)\) fails to exist at every single point of \([0,1]\) (the values \(f_k(x)\) are unbounded in \(k\)). This shows that in 7.23 one really must pass to a subsequence: convergence in \(\|\cdot\|_p\) gives pointwise almost-everywhere convergence only along some subsequence. Of course 7.23 does apply here, and indeed the subsequence consisting of the terms \(f_{2^n - 1} = n\chi_{[0, 2^{-n}]}\) converges to \(0\) at every \(x \in (0,1]\).

Problem (7B.13)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space, \(1 \le p \le \infty\), \(f \in \mathcal{L}^p(\mu)\), and \(f_1, f_2, \ldots\) is a sequence in \(\mathcal{L}^p(\mu)\) such that \(\lim_{k \to \infty} \|f_k - f\|_p = 0\). Show that if \(g : X \to \mathbf{F}\) is a function such that \(\lim_{k \to \infty} f_k(x) = g(x)\) for almost every \(x \in X\), then \(f(x) = g(x)\) for almost every \(x \in X\).

Solution

Because \(f \in \mathcal{L}^p(\mu)\), the sequence \(f_1, f_2, \ldots\) lies in \(\mathcal{L}^p(\mu)\), and \(\lim_{k \to \infty}\|f_k - f\|_p = 0\), the result 7.23 applies: there is a subsequence \(f_{k_1}, f_{k_2}, \ldots\) (with \(k_1 < k_2 < \cdots\)) such that

\begin{equation*} \lim_{m \to \infty} f_{k_m}(x) = f(x) \quad \text{for almost every } x \in X . \end{equation*}

In other words, there exists \(E \in \mathcal{S}\) with \(\mu(E) = 0\) such that \(\lim_{m \to \infty} f_{k_m}(x) = f(x)\) for every \(x \in X \setminus E\).

By hypothesis there also exists \(F \in \mathcal{S}\) with \(\mu(F) = 0\) such that

\begin{equation*} \lim_{k \to \infty} f_k(x) = g(x) \quad \text{for every } x \in X \setminus F . \end{equation*}

Let \(N = E \cup F\). Then \(N \in \mathcal{S}\) and, by countable subadditivity of \(\mu\) (2.58),

\begin{equation*} \mu(N) \le \mu(E) + \mu(F) = 0 , \end{equation*}

so \(\mu(N) = 0\).

Now suppose \(x \in X \setminus N\). Since \(x \notin F\), the sequence \(f_1(x), f_2(x), \ldots\) converges in \(\mathbf{F}\) to \(g(x)\). Every subsequence of a convergent sequence converges to the same limit, so

\begin{equation*} \lim_{m \to \infty} f_{k_m}(x) = g(x) . \end{equation*}

Since \(x \notin E\), we also have \(\lim_{m \to \infty} f_{k_m}(x) = f(x)\). Limits of sequences in \(\mathbf{F}\) are unique, so

\begin{equation*} f(x) = g(x) . \end{equation*}

Thus \(f(x) = g(x)\) for every \(x \in X \setminus N\), where \(\mu(N) = 0\). Hence \(f(x) = g(x)\) for almost every \(x \in X\), as required.

Remark. No measurability assumption on \(g\) was needed or used. The conclusion says that although convergence in \(\|\cdot\|_p\) need not give pointwise convergence anywhere (see Exercise 12), whenever the sequence does happen to converge pointwise almost everywhere, the pointwise limit must agree almost everywhere with the \(\mathcal{L}^p\)-limit. So the two notions of limit cannot disagree on a set of positive measure.

Problem (7B.14)

(a) Give an example of a measure \(\mu\) such that 7.25 fails for \(p = 1\).

(b) Show that if \(\mu\) is a \(\sigma\)-finite measure, then 7.25 holds for \(p = 1\).

Solution

Recall the statement of 7.25, now with \(p = 1\) (so that \(p^{\prime} = \infty\)): for \(h \in L^\infty(\mu)\) define \(\varphi_h : L^1(\mu) \to \mathbf{F}\) by \(\varphi_h(f) = \int f h \, d\mu\); the assertion is that \(h \mapsto \varphi_h\) is a one-to-one linear map from \(L^\infty(\mu)\) to \(\bigl(L^1(\mu)\bigr)^{\prime}\) with \(\|\varphi_h\| = \|h\|_\infty\) for every \(h\).

Two things are true for every measure \(\mu\) and require no hypotheses. First, Hölder’s inequality (7.9) with \(p = 1\) and \(p^{\prime} = \infty\) gives \(\|fh\|_1 \le \|f\|_1 \|h\|_\infty\) for all \(\mathcal{S}\)-measurable \(f, h\); hence \(fh \in \mathcal{L}^1(\mu)\) whenever \(f \in \mathcal{L}^1(\mu)\) and \(h \in \mathcal{L}^\infty(\mu)\), the map \(\varphi_h\) is a well-defined bounded linear functional on \(L^1(\mu)\), and

\begin{equation*} \|\varphi_h\| \le \|h\|_\infty . \end{equation*}

(The integral \(\int fh \, d\mu\) is unchanged if \(f\) or \(h\) is altered on a set of measure \(0\), so \(\varphi_h\) is well defined on \(L^1(\mu)\) and depends only on the class of \(h\) in \(L^\infty(\mu)\).) Second, \(h \mapsto \varphi_h\) is clearly linear, by linearity of the integral. What can fail is the reverse inequality \(\|h\|_\infty \le \|\varphi_h\|\), and with it the injectivity of \(h \mapsto \varphi_h\).

Part (a). Let \(X = \{1\}\), let \(\mathcal{S} = \{\emptyset, X\}\), and define \(\mu\) on \(\mathcal{S}\) by \(\mu(\emptyset) = 0\) and \(\mu(X) = \infty\). Then \(\mu\) is a measure: in any disjoint sequence of sets in \(\mathcal{S}\) at most one set is \(X\) and the rest are \(\emptyset\), so countable additivity holds.

Every function \(f : X \to \mathbf{F}\) is \(\mathcal{S}\)-measurable, and \(f\) is determined by the single value \(c = f(1)\). Then \(|f| = |c| \chi_X\), so

\begin{equation*} \int |f| \, d\mu = |c| \, \mu(X) = \begin{cases} 0 & \text{if } c = 0, \\ \infty & \text{if } c \ne 0, \end{cases} \end{equation*}

using the convention \(0 \cdot \infty = 0\). Hence

\begin{equation*} \mathcal{L}^1(\mu) = \{0\} , \qquad \text{and so} \qquad L^1(\mu) = \{0\} \quad \text{and} \quad \bigl(L^1(\mu)\bigr)^{\prime} = \{0\} . \end{equation*}

On the other hand, for \(f\) with \(f(1) = c\) we have \(\{x \in X : |f(x)| > t\}\) equal to \(X\) if \(t < |c|\) and equal to \(\emptyset\) if \(t \ge |c|\), so

\begin{equation*} \|f\|_\infty = \inf\{t > 0 : \mu(\{|f| > t\}) = 0\} = |c| . \end{equation*}

Thus \(\mathcal{L}^\infty(\mu)\) consists of all such \(f\), and since the only set of \(\mu\)-measure \(0\) is \(\emptyset\), distinct functions give distinct elements of \(L^\infty(\mu)\). So \(L^\infty(\mu)\) is a one-dimensional normed vector space, isometrically \(\mathbf{F}\) with its usual absolute value.

Now take \(h\) with \(h(1) = 1\), so \(\|h\|_\infty = 1\). Since \(L^1(\mu) = \{0\}\), the only linear functional on \(L^1(\mu)\) is \(0\), so \(\varphi_h = 0\) and

\begin{equation*} \|\varphi_h\| = 0 \ne 1 = \|h\|_\infty . \end{equation*}

So the norm-preserving assertion of 7.25 fails. Moreover \(\varphi_h = \varphi_0\) although \(h \ne 0\) in \(L^\infty(\mu)\), so \(h \mapsto \varphi_h\) is not one-to-one either. Hence 7.25 fails for \(p = 1\) for this measure \(\mu\).

(The mechanism is that \(\mu\) is not semifinite: \(X\) has infinite measure but contains no measurable subset of positive finite measure, so \(L^1(\mu)\) is too small to detect \(h\).)

Part (b). Suppose now that \(\mu\) is \(\sigma\)-finite, say \(X = \bigcup_{n=1}^{\infty} X_n\) with \(X_n \in \mathcal{S}\) and \(\mu(X_n) < \infty\) for each \(n\). Let \(h \in \mathcal{L}^\infty(\mu)\). By the first paragraph it suffices to prove \(\|h\|_\infty \le \|\varphi_h\|\).

If \(\|h\|_\infty = 0\) this is immediate, so assume \(\|h\|_\infty > 0\) and let \(t\) satisfy \(0 < t < \|h\|_\infty\). Set

\begin{equation*} E = \{x \in X : |h(x)| > t\} \in \mathcal{S} . \end{equation*}

If we had \(\mu(E) = 0\), then the definition of the essential supremum (7.1) would give \(\|h\|_\infty \le t\), contradicting \(t < \|h\|_\infty\). Hence \(\mu(E) > 0\). Because \(E = \bigcup_{n=1}^{\infty} (E \cap X_n)\), countable subadditivity (2.58) shows that \(\mu(E \cap X_n) > 0\) for some \(n\). Fix such an \(n\) and set \(A = E \cap X_n\), so that

\begin{equation*} A \in \mathcal{S} , \qquad 0 < \mu(A) \le \mu(X_n) < \infty . \end{equation*}

Define \(f : X \to \mathbf{F}\) by

\begin{equation*} f(x) = \begin{cases} \dfrac{1}{\mu(A)} \cdot \dfrac{\overline{h(x)}}{|h(x)|} & \text{if } x \in A, \\[2mm] 0 & \text{if } x \in X \setminus A. \end{cases} \end{equation*}

This is well defined because \(|h(x)| > t > 0\) for every \(x \in A\), and \(f\) is \(\mathcal{S}\)-measurable because \(h\) is \(\mathcal{S}\)-measurable and \(A \in \mathcal{S}\). (If \(\mathbf{F} = \mathbf{R}\), read \(\overline{h(x)}/|h(x)|\) as the sign of \(h(x)\).) Since \(|f| = \frac{1}{\mu(A)} \chi_A\), we get

\begin{equation*} \|f\|_1 = \frac{\mu(A)}{\mu(A)} = 1 , \end{equation*}

so \(f \in \mathcal{L}^1(\mu)\) with \(\|f\|_1 = 1\). Moreover \(f(x) h(x) = \frac{|h(x)|}{\mu(A)}\) for \(x \in A\) and \(f(x)h(x) = 0\) otherwise, so

\begin{equation*} \varphi_h(f) = \int f h \, d\mu = \frac{1}{\mu(A)} \int_A |h| \, d\mu \ge \frac{1}{\mu(A)} \cdot t \, \mu(A) = t , \end{equation*}

where the inequality uses \(|h| > t\) on \(A\) together with \(\mu(A) < \infty\). Therefore

\begin{equation*} \|\varphi_h\| \ge \frac{|\varphi_h(f)|}{\|f\|_1} \ge t . \end{equation*}

Since this holds for every \(t \in (0, \|h\|_\infty)\), we conclude \(\|\varphi_h\| \ge \|h\|_\infty\). Combined with \(\|\varphi_h\| \le \|h\|_\infty\) from Hölder’s inequality, this gives

\begin{equation*} \|\varphi_h\| = \|h\|_\infty \qquad \text{for every } h \in L^\infty(\mu) . \end{equation*}

Finally, \(h \mapsto \varphi_h\) is one-to-one: if \(h_1, h_2 \in L^\infty(\mu)\) satisfy \(\varphi_{h_1} = \varphi_{h_2}\), then by linearity

\begin{equation*} \|h_1 - h_2\|_\infty = \|\varphi_{h_1 - h_2}\| = \|\varphi_{h_1} - \varphi_{h_2}\| = \|0\| = 0 , \end{equation*}

which forces \(h_1 = h_2\) as elements of \(L^\infty(\mu)\). Thus \(h \mapsto \varphi_h\) is a one-to-one linear map from \(L^\infty(\mu)\) into \(\bigl(L^1(\mu)\bigr)^{\prime}\) that preserves norms, which is precisely 7.25 for \(p = 1\).

Problem (7B.15)

Let

\begin{equation*} c_0 = \{(a_1, a_2, \ldots) \in \ell^\infty : \lim_{k \to \infty} a_k = 0\}. \end{equation*}

Give \(c_0\) the norm that it inherits as a subspace of \(\ell^\infty\).

(a) Prove that \(c_0\) is a Banach space.

(b) Prove that the dual space of \(c_0\) can be identified with \(\ell^1\).

Solution

Throughout, \(\ell^\infty\) carries the norm \(\|a\|_\infty = \sup_k |a_k|\) and \(\ell^1\) carries \(\|b\|_1 = \sum_{k=1}^\infty |b_k|\). Recall that \(\ell^\infty = L^\infty(\mu)\) and \(\ell^1 = L^1(\mu)\), where \(\mu\) is counting measure on \(\mathbb{Z}^+\) (no nonempty set has measure \(0\), so the essential supremum is the supremum); thus \(\ell^\infty\) and \(\ell^1\) are Banach spaces by 7.24.

Part (a). First, \(c_0\) is a subspace of \(\ell^\infty\): the sequence \((0,0,\ldots)\) lies in \(c_0\), and if \(a, a^{\prime} \in c_0\) and \(\alpha \in \mathbf{F}\) then \(\lim_{k \to \infty}(a_k + \alpha a_k^{\prime}) = 0\), so \(a + \alpha a^{\prime} \in c_0\).

Next, \(c_0\) is a closed subset of \(\ell^\infty\). Suppose \(a^{(1)}, a^{(2)}, \ldots\) is a sequence in \(c_0\) and \(a \in \ell^\infty\) with \(\lim_{n \to \infty} \|a^{(n)} - a\|_\infty = 0\). Let \(\varepsilon > 0\). Choose \(n\) with \(\|a^{(n)} - a\|_\infty < \varepsilon/2\), and then choose \(K \in \mathbb{Z}^+\) such that \(|a^{(n)}_k| < \varepsilon/2\) for all \(k > K\) (possible because \(a^{(n)} \in c_0\)). For \(k > K\) we then have

\begin{equation*} |a_k| \le |a_k - a^{(n)}_k| + |a^{(n)}_k| < \tfrac{\varepsilon}{2} + \tfrac{\varepsilon}{2} = \varepsilon . \end{equation*}

Hence \(\lim_{k \to \infty} a_k = 0\), so \(a \in c_0\). Thus \(c_0\) is closed in \(\ell^\infty\).

Finally, a closed subspace of a Banach space is a Banach space. Indeed, suppose \(a^{(1)}, a^{(2)}, \ldots\) is a Cauchy sequence in \(c_0\). Because the norm on \(c_0\) is the restriction of the norm on \(\ell^\infty\), this sequence is Cauchy in \(\ell^\infty\), which is complete; hence it converges in \(\ell^\infty\) to some \(a \in \ell^\infty\). Because \(c_0\) is closed, \(a \in c_0\), and \(\lim_{n \to \infty}\|a^{(n)} - a\|_\infty = 0\) is convergence in \(c_0\). Thus \(c_0\) is complete, so \(c_0\) is a Banach space.

Part (b). For \(k \in \mathbb{Z}^+\) let \(e_k \in c_0\) be the sequence whose \(k^{\text{th}}\) term is \(1\) and whose other terms are \(0\).

For \(b = (b_1, b_2, \ldots) \in \ell^1\) define \(\varphi_b : c_0 \to \mathbf{F}\) by

\begin{equation*} \varphi_b(a) = \sum_{k=1}^\infty a_k b_k . \end{equation*}

The series converges absolutely, because \(\sum_{k=1}^\infty |a_k b_k| \le \|a\|_\infty \|b\|_1 < \infty\). The same estimate shows \(|\varphi_b(a)| \le \|b\|_1 \|a\|_\infty\); clearly \(\varphi_b\) is linear in \(a\). Hence \(\varphi_b \in (c_0)^{\prime}\) with \(\|\varphi_b\| \le \|b\|_1\). Also \(b \mapsto \varphi_b\) is linear.

To see that \(\|\varphi_b\| = \|b\|_1\), fix \(n \in \mathbb{Z}^+\) and define \(a^{(n)} \in c_0\) by

\begin{equation*} a^{(n)}_k = \begin{cases} \overline{b_k}/|b_k| & \text{if } k \le n \text{ and } b_k \ne 0,\\ 0 & \text{otherwise.}\end{cases} \end{equation*}

Only finitely many terms of \(a^{(n)}\) are nonzero, so \(a^{(n)} \in c_0\), and \(\|a^{(n)}\|_\infty \le 1\). Now

\begin{equation*} \sum_{k=1}^n |b_k| = \varphi_b\bigl(a^{(n)}\bigr) \le \|\varphi_b\| \, \|a^{(n)}\|_\infty \le \|\varphi_b\| . \end{equation*}

Letting \(n \to \infty\) gives \(\|b\|_1 \le \|\varphi_b\|\). Hence \(\|\varphi_b\| = \|b\|_1\) for all \(b \in \ell^1\). In particular \(b \mapsto \varphi_b\) is one-to-one (if \(\varphi_b = 0\) then \(\|b\|_1 = 0\), so \(b = 0\)).

It remains to prove that \(b \mapsto \varphi_b\) maps \(\ell^1\) onto \((c_0)^{\prime}\). Suppose \(\varphi \in (c_0)^{\prime}\). Define \(b = (b_1, b_2, \ldots)\) by

\begin{equation*} b_k = \varphi(e_k). \end{equation*}

With \(a^{(n)}\) defined from this \(b\) as above, \(a^{(n)}\) is a finite linear combination of \(e_1, \ldots, e_n\), so linearity of \(\varphi\) gives

\begin{equation*} \sum_{k=1}^n |b_k| = \varphi\bigl(a^{(n)}\bigr) \le \|\varphi\| \, \|a^{(n)}\|_\infty \le \|\varphi\| . \end{equation*}

Letting \(n \to \infty\) shows \(\|b\|_1 \le \|\varphi\| < \infty\), so \(b \in \ell^1\).

Now \(\varphi\) and \(\varphi_b\) are bounded linear functionals on \(c_0\) that agree at each \(e_k\) (both give \(b_k\)), hence they agree on the subspace of finitely supported sequences. That subspace is dense in \(c_0\): if \(a \in c_0\) and \(a^{[n]} = (a_1, \ldots, a_n, 0, 0, \ldots)\), then

\begin{equation*} \|a - a^{[n]}\|_\infty = \sup_{k > n} |a_k| \longrightarrow 0 \end{equation*}

as \(n \to \infty\), because \(\lim_{k \to \infty} a_k = 0\). Both \(\varphi\) and \(\varphi_b\) are bounded linear functionals, hence continuous by 6.48, and two continuous functions into \(\mathbf{F}\) that agree on a dense subset agree everywhere; so \(\varphi = \varphi_b\).

Thus \(b \mapsto \varphi_b\) is a linear bijection from \(\ell^1\) onto \((c_0)^{\prime}\) that preserves norms. In other words, the dual space of \(c_0\) can be identified with \(\ell^1\).

Problem (7B.16)

Suppose \(1 \le p \le 2\).

(a) Prove that if \(w, z \in \mathbb{C}\), then

\begin{equation*} \frac{|w + z|^p + |w - z|^p}{2} \le |w|^p + |z|^p \le \frac{|w + z|^p + |w - z|^p}{2^{p-1}} . \end{equation*}

(b) Prove that if \(\mu\) is a measure and \(f, g \in L^p(\mu)\), then

\begin{equation*} \frac{\|f + g\|_p^p + \|f - g\|_p^p}{2} \le \|f\|_p^p + \|g\|_p^p \le \frac{\|f + g\|_p^p + \|f - g\|_p^p}{2^{p-1}} . \end{equation*}

Solution

We first isolate two elementary facts.

Lemma 1. If \(s \in (0, 1]\) and \(\alpha, \beta \ge 0\), then \((\alpha + \beta)^s \le \alpha^s + \beta^s\); if \(s \ge 1\) and \(\alpha, \beta \ge 0\), then \(\alpha^s + \beta^s \le (\alpha + \beta)^s\).

Proof of Lemma 1. If \(\alpha + \beta = 0\) both statements are trivial, so assume \(\alpha + \beta > 0\) and set \(t = \alpha/(\alpha + \beta) \in [0,1]\), so that \(1 - t = \beta/(\alpha + \beta)\). If \(0 < s \le 1\), then \(t^s \ge t\) and \((1 - t)^s \ge 1 - t\) (raising a number in \([0,1]\) to a power at most \(1\) does not decrease it); adding gives \(t^s + (1-t)^s \ge 1\), and multiplying by \((\alpha + \beta)^s\) gives \(\alpha^s + \beta^s \ge (\alpha+\beta)^s\). If \(s \ge 1\), then \(t^s \le t\) and \((1-t)^s \le 1 - t\), so \(t^s + (1-t)^s \le 1\), and multiplying by \((\alpha+\beta)^s\) gives \(\alpha^s + \beta^s \le (\alpha+\beta)^s\).

Lemma 2. If \(0 < r \le 1\) and \(\alpha, \beta \ge 0\), then

\begin{equation*} \frac{\alpha^r + \beta^r}{2} \le \Bigl(\frac{\alpha + \beta}{2}\Bigr)^{r} . \end{equation*}

Proof of Lemma 2. The function \(t \mapsto t^r\) on \([0, \infty)\) is concave when \(0 < r \le 1\) (its second derivative \(r(r-1)t^{r-2}\) is at most \(0\) on \((0, \infty)\), and the function is continuous at \(0\)). Concavity applied to the two points \(\alpha, \beta\) with equal weights \(\tfrac12\) gives the stated inequality.

Part (a). Fix \(p\) with \(1 \le p \le 2\), and put \(r = p/2\), so \(0 < r \le 1\).

First we prove the left-hand inequality, which is equivalent to

\begin{equation*} |w + z|^p + |w - z|^p \le 2\bigl(|w|^p + |z|^p\bigr) . \end{equation*}

Apply Lemma 2 with \(\alpha = |w+z|^2\) and \(\beta = |w-z|^2\), and use the parallelogram equality \(|w+z|^2 + |w-z|^2 = 2\bigl(|w|^2 + |z|^2\bigr)\) (which holds in \(\mathbb{C}\), as a direct expansion shows):

\begin{equation*} \frac{|w+z|^p + |w-z|^p}{2} = \frac{\alpha^r + \beta^r}{2} \le \Bigl(\frac{\alpha + \beta}{2}\Bigr)^{r} = \bigl(|w|^2 + |z|^2\bigr)^{p/2} . \end{equation*}

Now apply Lemma 1 with \(s = p/2 \in (0,1]\), \(\alpha = |w|^2\), \(\beta = |z|^2\) to get

\begin{equation*} \bigl(|w|^2 + |z|^2\bigr)^{p/2} \le |w|^p + |z|^p . \end{equation*}

Combining the two displays gives \(\frac{|w+z|^p + |w-z|^p}{2} \le |w|^p + |z|^p\), as desired.

Now the right-hand inequality. Apply the inequality just proved with \(w\) replaced by \(w + z\) and \(z\) replaced by \(w - z\). Because \((w+z) + (w-z) = 2w\) and \((w+z) - (w-z) = 2z\), this yields

\begin{equation*} \frac{|2w|^p + |2z|^p}{2} \le |w+z|^p + |w-z|^p , \end{equation*}

that is, \(2^{p-1}\bigl(|w|^p + |z|^p\bigr) \le |w+z|^p + |w-z|^p\). Dividing by \(2^{p-1}\) gives

\begin{equation*} |w|^p + |z|^p \le \frac{|w+z|^p + |w-z|^p}{2^{p-1}} , \end{equation*}

completing the proof of (a). (If \(\mathbf{F} = \mathbb{R}\) the same inequalities hold, as \(\mathbb{R} \subseteq \mathbb{C}\).)

Part (b). Because \(1 \le p \le 2\) we have \(p < \infty\), so \(\|h\|_p^p = \int |h|^p \, d\mu\) for \(h \in L^p(\mu)\).

Let \((X, \mathcal{S}, \mu)\) be the measure space in question, and choose representatives of \(f\) and \(g\) defined on all of \(X\); changing a representative on a set of measure \(0\) changes none of the four quantities in the desired display, so this is harmless. By part (a) applied at each \(x \in X\) (with \(w = f(x)\) and \(z = g(x)\)),

\begin{equation*} |f(x) + g(x)|^p + |f(x) - g(x)|^p \le 2\bigl(|f(x)|^p + |g(x)|^p\bigr) \end{equation*}

and

\begin{equation*} 2^{p-1}\bigl(|f(x)|^p + |g(x)|^p\bigr) \le |f(x) + g(x)|^p + |f(x) - g(x)|^p . \end{equation*}

All four of the functions \(|f+g|^p\), \(|f-g|^p\), \(|f|^p\), \(|g|^p\) are nonnegative measurable functions with finite integral (the first two because \(f + g\) and \(f - g\) belong to \(L^p(\mu)\), since \(L^p(\mu)\) is a vector space by 7.18). The pointwise inequalities above hold at every \(x \in X\) for the chosen representatives, so integrating them with respect to \(\mu\) and using that integration is order preserving (3.8) together with additivity of integration for nonnegative measurable functions (3.16) gives

\begin{equation*} \|f+g\|_p^p + \|f-g\|_p^p \le 2\bigl(\|f\|_p^p + \|g\|_p^p\bigr) \end{equation*}

and

\begin{equation*} 2^{p-1}\bigl(\|f\|_p^p + \|g\|_p^p\bigr) \le \|f+g\|_p^p + \|f-g\|_p^p . \end{equation*}

Dividing the first by \(2\) and the second by \(2^{p-1}\) gives exactly

\begin{equation*} \frac{\|f + g\|_p^p + \|f - g\|_p^p}{2} \le \|f\|_p^p + \|g\|_p^p \le \frac{\|f + g\|_p^p + \|f - g\|_p^p}{2^{p-1}} , \end{equation*}

as desired.

Problem (7B.17)

Suppose \(2 \le p < \infty\).

(a) Prove that if \(w, z \in \mathbb{C}\), then

\begin{equation*} \frac{|w + z|^p + |w - z|^p}{2^{p-1}} \le |w|^p + |z|^p \le \frac{|w + z|^p + |w - z|^p}{2} . \end{equation*}

(b) Prove that if \(\mu\) is a measure and \(f, g \in L^p(\mu)\), then

\begin{equation*} \frac{\|f + g\|_p^p + \|f - g\|_p^p}{2^{p-1}} \le \|f\|_p^p + \|g\|_p^p \le \frac{\|f + g\|_p^p + \|f - g\|_p^p}{2} . \end{equation*}

[The inequalities in the two previous exercises are called Clarkson’s inequalities. They were discovered by James Clarkson in 1936.]

Solution

As in the previous exercise we begin with two elementary facts.

Lemma 1. If \(s \ge 1\) and \(\alpha, \beta \ge 0\), then \(\alpha^s + \beta^s \le (\alpha + \beta)^s\).

Proof of Lemma 1. If \(\alpha + \beta = 0\) this is trivial. Otherwise set \(t = \alpha/(\alpha+\beta) \in [0,1]\), so \(1 - t = \beta/(\alpha+\beta)\). Because \(s \ge 1\) and \(t, 1-t \in [0,1]\), we have \(t^s \le t\) and \((1-t)^s \le 1-t\); adding gives \(t^s + (1-t)^s \le 1\), and multiplying by \((\alpha+\beta)^s\) gives \(\alpha^s + \beta^s \le (\alpha+\beta)^s\).

Lemma 2. If \(r \ge 1\) and \(\alpha, \beta \ge 0\), then

\begin{equation*} \Bigl(\frac{\alpha + \beta}{2}\Bigr)^{r} \le \frac{\alpha^r + \beta^r}{2} . \end{equation*}

Proof of Lemma 2. The function \(t \mapsto t^r\) on \([0,\infty)\) is convex when \(r \ge 1\) (its second derivative \(r(r-1)t^{r-2}\) is nonnegative on \((0, \infty)\), and the function is continuous at \(0\)). Applying convexity to the points \(\alpha, \beta\) with equal weights \(\tfrac12\) gives the stated inequality.

Part (a). Fix \(p\) with \(2 \le p < \infty\) and put \(r = p/2 \ge 1\).

We first prove the right-hand inequality, which is equivalent to

\begin{equation*} 2\bigl(|w|^p + |z|^p\bigr) \le |w+z|^p + |w-z|^p . \end{equation*}

Set \(\alpha = |w+z|^2\) and \(\beta = |w-z|^2\). The parallelogram equality (verified by expanding \(|u|^2 = u \overline{u}\)) gives \(\alpha + \beta = 2\bigl(|w|^2 + |z|^2\bigr)\). Lemma 2 therefore yields

\begin{equation*} \frac{|w+z|^p + |w-z|^p}{2} = \frac{\alpha^r + \beta^r}{2} \ge \Bigl(\frac{\alpha+\beta}{2}\Bigr)^{r} = \bigl(|w|^2 + |z|^2\bigr)^{p/2} . \end{equation*}

Lemma 1 with \(s = p/2 \ge 1\), \(\alpha = |w|^2\), \(\beta = |z|^2\) gives

\begin{equation*} \bigl(|w|^2 + |z|^2\bigr)^{p/2} \ge |w|^p + |z|^p . \end{equation*}

Combining the two displays gives \(|w|^p + |z|^p \le \frac{|w+z|^p + |w-z|^p}{2}\), which is the right-hand inequality.

For the left-hand inequality, apply the inequality just proved with \(w\) replaced by \(w+z\) and \(z\) replaced by \(w-z\). Since \((w+z)+(w-z) = 2w\) and \((w+z)-(w-z) = 2z\), we obtain

\begin{equation*} |w+z|^p + |w-z|^p \le \frac{|2w|^p + |2z|^p}{2} = 2^{p-1}\bigl(|w|^p + |z|^p\bigr) . \end{equation*}

Dividing by \(2^{p-1}\) gives

\begin{equation*} \frac{|w+z|^p + |w-z|^p}{2^{p-1}} \le |w|^p + |z|^p , \end{equation*}

completing the proof of (a). (The same inequalities hold for \(w, z \in \mathbb{R}\), because \(\mathbb{R} \subseteq \mathbb{C}\).)

Part (b). Because \(p < \infty\) we have \(\|h\|_p^p = \int |h|^p \, d\mu\) for \(h \in L^p(\mu)\), and \(f + g, f - g \in L^p(\mu)\) because \(L^p(\mu)\) is a vector space by 7.18.

Let \((X, \mathcal{S}, \mu)\) be the measure space in question, and choose representatives of \(f\) and \(g\) defined on all of \(X\); changing a representative on a set of measure \(0\) changes none of the four quantities in the desired display, so this is harmless. Applying part (a) with \(w = f(x)\) and \(z = g(x)\) gives, for every \(x \in X\),

\begin{equation*} |f(x)+g(x)|^p + |f(x)-g(x)|^p \le 2^{p-1}\bigl(|f(x)|^p + |g(x)|^p\bigr) \end{equation*}

and

\begin{equation*} 2\bigl(|f(x)|^p + |g(x)|^p\bigr) \le |f(x)+g(x)|^p + |f(x)-g(x)|^p . \end{equation*}

Each of \(|f+g|^p, |f-g|^p, |f|^p, |g|^p\) is a nonnegative \(\mathcal{S}\)-measurable function with finite integral. The pointwise inequalities above hold at every \(x \in X\) for the chosen representatives, so integrating them with respect to \(\mu\) and using that integration is order preserving (3.8) together with additivity of integration for nonnegative measurable functions (3.16) gives

\begin{equation*} \|f+g\|_p^p + \|f-g\|_p^p \le 2^{p-1}\bigl(\|f\|_p^p + \|g\|_p^p\bigr) \end{equation*}

and

\begin{equation*} 2\bigl(\|f\|_p^p + \|g\|_p^p\bigr) \le \|f+g\|_p^p + \|f-g\|_p^p . \end{equation*}

Dividing the first by \(2^{p-1}\) and the second by \(2\) produces

\begin{equation*} \frac{\|f + g\|_p^p + \|f - g\|_p^p}{2^{p-1}} \le \|f\|_p^p + \|g\|_p^p \le \frac{\|f + g\|_p^p + \|f - g\|_p^p}{2} , \end{equation*}

as desired.

Note that when \(p = 2\) all four inequalities in this exercise and the previous one become the parallelogram equality, since then \(2^{p-1} = 2\).

Problem (7B.18)

Suppose \((X, \mathcal{S}, \mu)\) is a measure space, \(1 \le p, q \le \infty\), and \(h : X \to \mathbf{F}\) is an \(\mathcal{S}\)-measurable function such that \(hf \in L^q(\mu)\) for every \(f \in L^p(\mu)\). Prove that \(f \mapsto hf\) is a continuous linear map from \(L^p(\mu)\) to \(L^q(\mu)\).

Solution

Define \(T : L^p(\mu) \to L^q(\mu)\) by \(Tf = hf\); the hypothesis says exactly that \(T\) maps \(L^p(\mu)\) into \(L^q(\mu)\).

\(T\) is well defined on \(L^p(\mu)\) regarded as a space of equivalence classes: if \(f_1 = f_2\) almost everywhere, then \(hf_1 = hf_2\) almost everywhere, so \(hf_1\) and \(hf_2\) determine the same element of \(L^q(\mu)\). Furthermore \(T\) is linear, because

\begin{equation*} h \cdot (f_1 + \alpha f_2) = hf_1 + \alpha\, hf_2 \end{equation*}

pointwise, for all \(f_1, f_2 \in L^p(\mu)\) and \(\alpha \in \mathbf{F}\).

By 6.48, to prove that \(T\) is continuous it suffices to prove that \(T\) is bounded. Both \(L^p(\mu)\) and \(L^q(\mu)\) are Banach spaces by 7.24, so the Closed Graph Theorem (6.85) applies: \(T\) is bounded if and only if

\begin{equation*} \operatorname{graph}(T) = \{(f, Tf) : f \in L^p(\mu)\} \end{equation*}

is a closed subspace of \(L^p(\mu) \times L^q(\mu)\). Because \(T\) is linear, \(\operatorname{graph}(T)\) is a subspace. So we need only prove that it is closed.

Suppose \((f_1, Tf_1), (f_2, Tf_2), \ldots\) is a sequence in \(\operatorname{graph}(T)\) converging to some \((f, g) \in L^p(\mu) \times L^q(\mu)\). By 6.84 this means

\begin{equation*} \lim_{k \to \infty} \|f_k - f\|_p = 0 \qquad \text{and} \qquad \lim_{k \to \infty} \|hf_k - g\|_q = 0 . \end{equation*}

We must show \(g = hf\) (as elements of \(L^q(\mu)\), i.e. almost everywhere).

Since \(\lim_{k \to \infty}\|f_k - f\|_p = 0\), by 7.23 there is a subsequence \(f_{k_1}, f_{k_2}, \ldots\) such that

\begin{equation*} \lim_{m \to \infty} f_{k_m}(x) = f(x) \quad \text{for almost every } x \in X . \end{equation*}

Multiplying by \(h(x)\) gives

\begin{equation*} \lim_{m \to \infty} h(x) f_{k_m}(x) = h(x) f(x) \quad \text{for almost every } x \in X . \end{equation*}

The sequence \(hf_{k_1}, hf_{k_2}, \ldots\) is a subsequence of a sequence converging to \(g\) in \(L^q(\mu)\), so \(\lim_{m \to \infty} \|hf_{k_m} - g\|_q = 0\). Applying 7.23 once more (now in \(L^q(\mu)\)), there is a further subsequence \(hf_{k_{m_1}}, hf_{k_{m_2}}, \ldots\) such that

\begin{equation*} \lim_{j \to \infty} h(x) f_{k_{m_j}}(x) = g(x) \quad \text{for almost every } x \in X . \end{equation*}

Let \(E_1\) be the set of measure \(0\) off which \(\lim_{m \to \infty} h(x)f_{k_m}(x) = h(x)f(x)\) holds, let \(E_2\) be the set of measure \(0\) off which \(\lim_{j \to \infty} h(x)f_{k_{m_j}}(x) = g(x)\) holds, and put \(E = E_1 \cup E_2\). Then \(\mu(E) = 0\), because a union of two sets of measure \(0\) has measure \(0\). For \(x \in X \setminus E\), the sequence \(j \mapsto h(x)f_{k_{m_j}}(x)\) converges to \(h(x)f(x)\) (it is a subsequence of a sequence converging to \(h(x)f(x)\)) and also converges to \(g(x)\). By uniqueness of limits in \(\mathbf{F}\),

\begin{equation*} g(x) = h(x) f(x) \quad \text{for all } x \in X \setminus E . \end{equation*}

Thus \(g = hf\) almost everywhere, so \(g = Tf\) in \(L^q(\mu)\) and \((f, g) = (f, Tf) \in \operatorname{graph}(T)\).

Hence \(\operatorname{graph}(T)\) is a closed subspace of \(L^p(\mu) \times L^q(\mu)\). The Closed Graph Theorem (6.85) now implies that \(T\) is a bounded linear map, and 6.48 implies that \(T\) is continuous. In other words, there exists \(c < \infty\) such that

\begin{equation*} \|hf\|_q \le c \|f\|_p \quad \text{for all } f \in L^p(\mu), \end{equation*}

and \(f \mapsto hf\) is a continuous linear map from \(L^p(\mu)\) to \(L^q(\mu)\).

Problem (7B.19)

Prove that if \(1 < p < \infty\), then \(\ell^p\) is reflexive.

[A Banach space is called reflexive if the canonical isometry of the Banach space into its double dual space is surjective (see Exercise 20 in Section 6D for the definitions of the double dual space and the canonical isometry).]

Solution

Fix \(p\) with \(1 < p < \infty\) and let \(p^{\prime} = p/(p-1)\) be its Hölder conjugate; then \(1 < p^{\prime} < \infty\) and \((p^{\prime})^{\prime} = p\).

For \(b \in \ell^{p^{\prime}}\) define \(\varphi_b : \ell^p \to \mathbf{F}\) by \(\varphi_b(a) = \sum_{k=1}^\infty a_k b_k\), and let

\begin{equation*} \Lambda_p : \ell^{p^{\prime}} \to (\ell^p)^{\prime}, \qquad \Lambda_p b = \varphi_b . \end{equation*}

By 7.26, \(\Lambda_p\) is a one-to-one linear map from \(\ell^{p^{\prime}}\) onto \((\ell^p)^{\prime}\) with \(\|\Lambda_p b\| = \|b\|_{p^{\prime}}\) for every \(b \in \ell^{p^{\prime}}\). Because \(1 \le p^{\prime} < \infty\) and \((p^{\prime})^{\prime} = p\), we may also apply 7.26 with \(p^{\prime}\) in place of \(p\): the map

\begin{equation*} \Lambda_{p^{\prime}} : \ell^{p} \to (\ell^{p^{\prime}})^{\prime}, \qquad (\Lambda_{p^{\prime}} a)(b) = \sum_{k=1}^\infty b_k a_k \quad \text{for } b \in \ell^{p^{\prime}}, \end{equation*}

is a one-to-one linear map from \(\ell^p\) onto \((\ell^{p^{\prime}})^{\prime}\) with \(\|\Lambda_{p^{\prime}} a\| = \|a\|_p\).

Let \(\Phi : \ell^p \to (\ell^p)^{\prime\prime}\) denote the canonical isometry, defined by

\begin{equation*} (\Phi a)(\varphi) = \varphi(a) \quad \text{for } a \in \ell^p \text{ and } \varphi \in (\ell^p)^{\prime} . \end{equation*}

By Exercise 20 in Section 6D, \(\|\Phi a\| = \|a\|_p\) for every \(a \in \ell^p\); in particular \(\Phi\) is one-to-one. Thus to prove that \(\ell^p\) is reflexive we must show that \(\Phi\) is surjective.

Suppose \(\psi \in (\ell^p)^{\prime\prime}\). Consider

\begin{equation*} \psi \circ \Lambda_p : \ell^{p^{\prime}} \to \mathbf{F} . \end{equation*}

This map is linear, being a composition of linear maps, and it is bounded because

\begin{equation*} |\psi(\Lambda_p b)| \le \|\psi\| \, \|\Lambda_p b\| = \|\psi\| \, \|b\|_{p^{\prime}} \end{equation*}

for all \(b \in \ell^{p^{\prime}}\). Hence \(\psi \circ \Lambda_p \in (\ell^{p^{\prime}})^{\prime}\).

Because \(\Lambda_{p^{\prime}}\) maps \(\ell^p\) onto \((\ell^{p^{\prime}})^{\prime}\), there exists \(a \in \ell^p\) with

\begin{equation*} \Lambda_{p^{\prime}} a = \psi \circ \Lambda_p . \end{equation*}

Evaluating both sides at an arbitrary \(b \in \ell^{p^{\prime}}\) gives

\begin{equation*} \sum_{k=1}^\infty b_k a_k = \psi(\varphi_b) . \end{equation*}

On the other hand, by the definition of \(\Phi\),

\begin{equation*} (\Phi a)(\varphi_b) = \varphi_b(a) = \sum_{k=1}^\infty a_k b_k . \end{equation*}

The two right-hand sides are equal, so

\begin{equation*} (\Phi a)(\varphi_b) = \psi(\varphi_b) \quad \text{for every } b \in \ell^{p^{\prime}} . \end{equation*}

Because \(\Lambda_p\) maps \(\ell^{p^{\prime}}\) onto \((\ell^p)^{\prime}\), every \(\varphi \in (\ell^p)^{\prime}\) equals \(\varphi_b\) for some \(b \in \ell^{p^{\prime}}\). Hence \((\Phi a)(\varphi) = \psi(\varphi)\) for every \(\varphi \in (\ell^p)^{\prime}\), which means \(\Phi a = \psi\).

Thus \(\Phi\) is surjective. Because \(\Phi\) is a norm-preserving linear map from \(\ell^p\) onto \((\ell^p)^{\prime\prime}\), we conclude that \(\ell^p\) is reflexive.

Problem (7B.20)

Prove that \(\ell^1\) is not reflexive.

Solution

For \(b \in \ell^\infty\) define \(\varphi_b : \ell^1 \to \mathbf{F}\) by \(\varphi_b(a) = \sum_{k=1}^\infty a_k b_k\), and let

\begin{equation*} \Lambda : \ell^\infty \to (\ell^1)^{\prime}, \qquad \Lambda b = \varphi_b . \end{equation*}

By 7.26 with \(p = 1\) (so that \(p^{\prime} = \infty\)), \(\Lambda\) is a one-to-one linear map from \(\ell^\infty\) onto \((\ell^1)^{\prime}\) with \(\|\Lambda b\| = \|b\|_\infty\) for all \(b \in \ell^\infty\). In particular \(\Lambda\) is invertible and \(\Lambda^{-1} : (\ell^1)^{\prime} \to \ell^\infty\) is a linear map with \(\|\Lambda^{-1}\varphi\|_\infty = \|\varphi\|\), hence bounded.

Let \(\Phi : \ell^1 \to (\ell^1)^{\prime\prime}\) be the canonical isometry, \((\Phi a)(\varphi) = \varphi(a)\). We show \(\Phi\) is not surjective.

Let \(c_0 \subseteq \ell^\infty\) be as in Exercise 15 and let \(\mathbf{1} = (1, 1, 1, \ldots) \in \ell^\infty\). In the proof of Exercise 15(a) we saw that \(c_0\) is a closed subspace of \(\ell^\infty\); thus \(\overline{c_0} = c_0\). Because the terms of \(\mathbf{1}\) do not have limit \(0\), we have \(\mathbf{1} \notin c_0 = \overline{c_0}\). Hence 6.73, applied with \(V = \ell^\infty\) and \(U = c_0\), provides

\begin{equation*} \tau \in (\ell^\infty)^{\prime} \quad \text{with} \quad \tau|_{c_0} = 0 \ \text{ and } \ \tau(\mathbf{1}) \ne 0 . \end{equation*}

Now define \(\psi : (\ell^1)^{\prime} \to \mathbf{F}\) by

\begin{equation*} \psi(\varphi) = \tau\bigl(\Lambda^{-1}\varphi\bigr) . \end{equation*}

Then \(\psi\) is linear (a composition of linear maps) and bounded, because

\begin{equation*} |\psi(\varphi)| \le \|\tau\| \, \|\Lambda^{-1}\varphi\|_\infty = \|\tau\| \, \|\varphi\| \end{equation*}

for all \(\varphi \in (\ell^1)^{\prime}\). Hence \(\psi \in (\ell^1)^{\prime\prime}\). Equivalently, \(\psi(\varphi_b) = \tau(b)\) for every \(b \in \ell^\infty\).

Suppose, for contradiction, that \(\Phi\) is surjective. Then there exists \(a = (a_1, a_2, \ldots) \in \ell^1\) with \(\Phi a = \psi\). Evaluating at \(\varphi_b\) for \(b \in \ell^\infty\) gives

\begin{equation*} \sum_{k=1}^\infty a_k b_k = \varphi_b(a) = (\Phi a)(\varphi_b) = \psi(\varphi_b) = \tau(b) . \end{equation*}

For \(j \in \mathbb{Z}^+\) let \(e_j \in \ell^\infty\) be the sequence whose \(j^{\text{th}}\) term is \(1\) and whose other terms are \(0\). Then \(e_j \in c_0\), so \(\tau(e_j) = 0\); taking \(b = e_j\) in the display above gives

\begin{equation*} a_j = \tau(e_j) = 0 . \end{equation*}

This holds for every \(j \in \mathbb{Z}^+\), so \(a = 0\). But then the display above gives \(\tau(b) = 0\) for every \(b \in \ell^\infty\), contradicting \(\tau(\mathbf{1}) \ne 0\).

This contradiction shows that no \(a \in \ell^1\) satisfies \(\Phi a = \psi\). Hence the canonical isometry \(\Phi\) of \(\ell^1\) into \((\ell^1)^{\prime\prime}\) is not surjective, which means that \(\ell^1\) is not reflexive.

Problem (7B.21)

Show that with the natural identifications, the canonical isometry of \(c_0\) into its double dual space is the inclusion map of \(c_0\) into \(\ell^\infty\) (see Exercise 15 for the definition of \(c_0\) and an identification of its dual space).

Solution

We first record the two natural identifications.

Identification of \((c_0)^{\prime}\) with \(\ell^1\). For \(b \in \ell^1\) define \(\psi_b : c_0 \to \mathbf{F}\) by \(\psi_b(a) = \sum_{k=1}^\infty a_k b_k\), and set

\begin{equation*} \Lambda_1 : \ell^1 \to (c_0)^{\prime}, \qquad \Lambda_1 b = \psi_b . \end{equation*}

By Exercise 15(b), \(\Lambda_1\) is a one-to-one linear map from \(\ell^1\) onto \((c_0)^{\prime}\) with \(\|\Lambda_1 b\| = \|b\|_1\) for every \(b \in \ell^1\).

Identification of \((\ell^1)^{\prime}\) with \(\ell^\infty\). For \(d \in \ell^\infty\) define \(\varphi_d : \ell^1 \to \mathbf{F}\) by \(\varphi_d(b) = \sum_{k=1}^\infty b_k d_k\), and set

\begin{equation*} \Lambda_\infty : \ell^\infty \to (\ell^1)^{\prime}, \qquad \Lambda_\infty d = \varphi_d . \end{equation*}

By 7.26 with \(p = 1\), \(\Lambda_\infty\) is a one-to-one linear map from \(\ell^\infty\) onto \((\ell^1)^{\prime}\) with \(\|\Lambda_\infty d\| = \|d\|_\infty\).

Identification of \((c_0)^{\prime\prime}\) with \((\ell^1)^{\prime}\). Define

\begin{equation*} \Gamma : (c_0)^{\prime\prime} \to (\ell^1)^{\prime}, \qquad \Gamma\theta = \theta \circ \Lambda_1 . \end{equation*}

Since \(\Lambda_1\) is linear and bounded and \(\theta\) is linear and bounded, \(\Gamma\theta\) is a bounded linear functional on \(\ell^1\); and \(\Gamma\) is clearly linear. Because \(\Lambda_1\) maps the closed unit ball of \(\ell^1\) onto the closed unit ball of \((c_0)^{\prime}\) (it is a norm-preserving bijection),

\begin{equation*} \|\Gamma\theta\| = \sup_{\|b\|_1 \le 1} |\theta(\Lambda_1 b)| = \sup_{\substack{\psi \in (c_0)^{\prime} \\ \|\psi\| \le 1}} |\theta(\psi)| = \|\theta\| . \end{equation*}

Moreover \(\Gamma\) is surjective: if \(\sigma \in (\ell^1)^{\prime}\), then \(\sigma \circ \Lambda_1^{-1} \in (c_0)^{\prime\prime}\) and \(\Gamma(\sigma \circ \Lambda_1^{-1}) = \sigma\). Hence \(\Gamma\) is a norm-preserving linear bijection from \((c_0)^{\prime\prime}\) onto \((\ell^1)^{\prime}\).

Combining the last two identifications, \((c_0)^{\prime\prime}\) is identified with \(\ell^\infty\) by the norm-preserving linear bijection

\begin{equation*} \Lambda_\infty^{-1} \circ \Gamma : (c_0)^{\prime\prime} \to \ell^\infty . \end{equation*}

Now let \(\Phi : c_0 \to (c_0)^{\prime\prime}\) be the canonical isometry, defined by \((\Phi a)(\psi) = \psi(a)\) for \(a \in c_0\) and \(\psi \in (c_0)^{\prime}\). Fix \(a \in c_0\). For every \(b \in \ell^1\) we have

\begin{equation*} \bigl(\Gamma(\Phi a)\bigr)(b) = (\Phi a)(\Lambda_1 b) = (\Phi a)(\psi_b) = \psi_b(a) = \sum_{k=1}^\infty a_k b_k . \end{equation*}

Since \(a \in c_0 \subseteq \ell^\infty\), the element \(\varphi_a = \Lambda_\infty a\) of \((\ell^1)^{\prime}\) is defined, and

\begin{equation*} (\Lambda_\infty a)(b) = \varphi_a(b) = \sum_{k=1}^\infty b_k a_k . \end{equation*}

The two right-hand sides agree for every \(b \in \ell^1\), so

\begin{equation*} \Gamma(\Phi a) = \Lambda_\infty a, \qquad \text{that is,} \qquad \bigl(\Lambda_\infty^{-1} \circ \Gamma\bigr)(\Phi a) = a . \end{equation*}

In other words, when \((c_0)^{\prime}\) is identified with \(\ell^1\) via \(\Lambda_1\) and \((c_0)^{\prime\prime} = ((c_0)^{\prime})^{\prime}\) is thereby identified with \((\ell^1)^{\prime}\) and hence with \(\ell^\infty\) via \(\Lambda_\infty\), the element \(\Phi a\) of \((c_0)^{\prime\prime}\) corresponds to the element \(a\) of \(\ell^\infty\). Thus under these natural identifications the canonical isometry \(\Phi\) becomes the map

\begin{equation*} c_0 \to \ell^\infty, \qquad a \mapsto a, \end{equation*}

which is exactly the inclusion map of \(c_0\) into \(\ell^\infty\).

In particular, the range of \(\Phi\) corresponds to \(c_0\), which is a proper subset of \(\ell^\infty\) (for example \((1,1,1,\ldots) \in \ell^\infty \setminus c_0\)). Hence \(\Phi\) is not surjective, so \(c_0\) is not reflexive.

Problem (7B.22)

Suppose \(1 \le p < \infty\) and \(V, W\) are Banach spaces. Show that \(V \times W\) is a Banach space if the norm on \(V \times W\) is defined by

\begin{equation*} \|(f,g)\| = \bigl(\|f\|^{p} + \|g\|^{p}\bigr)^{1/p} \end{equation*}

for \(f \in V\) and \(g \in W\).

Solution

Throughout, \(V \times W\) carries the natural vector space operations (Exercise 10 in Section 6B):

\begin{equation*} (f_1, g_1) + (f_2, g_2) = (f_1 + f_2,\, g_1 + g_2), \qquad \alpha (f,g) = (\alpha f, \alpha g) \end{equation*}

for \(f_1, f_2, f \in V\), \(g_1, g_2, g \in W\), and \(\alpha \in \mathbb{F}\). The additive identity is \((0,0)\) and the additive inverse of \((f,g)\) is \((-f,-g)\); verification of the vector space axioms is immediate from the corresponding axioms in \(V\) and in \(W\), since the operations act in each slot separately. So we must show two things: that the displayed formula defines a norm on this vector space, and that \(V \times W\) is complete with respect to it.

Two elementary inequalities. For \(a, b \in [0, \infty)\) we have

\begin{equation*} \max\{a,b\} \le \bigl(a^{p} + b^{p}\bigr)^{1/p} \le 2^{1/p} \max\{a,b\}. \end{equation*}

Indeed, if (say) \(a = \max\{a,b\}\) then \(a^{p} \le a^{p} + b^{p} \le 2a^{p}\), and raising to the power \(1/p\) preserves these inequalities because \(t \mapsto t^{1/p}\) is increasing on \([0,\infty)\). Writing \(|(f,g)|_{\infty} = \max\{\|f\|, \|g\|\}\) for the norm used in 6.84, the inequalities above say

\begin{equation*} |(f,g)|_{\infty} \le \|(f,g)\| \le 2^{1/p} \, |(f,g)|_{\infty} \end{equation*}

for all \((f,g) \in V \times W\); call this pair of inequalities \((\ast)\). Thus the two norms (once we know \(\|\cdot\|\) is a norm) are equivalent.

Positivity. Clearly \(\|(f,g)\| \in [0,\infty)\). If \(\|(f,g)\| = 0\) then \(\|f\|^{p} + \|g\|^{p} = 0\), and since both summands are nonnegative we get \(\|f\| = \|g\| = 0\), hence \(f = 0\) and \(g = 0\) because \(\|\cdot\|\) is a norm on \(V\) and on \(W\); thus \((f,g) = (0,0)\). Conversely \(\|(0,0)\| = 0\).

Homogeneity. For \(\alpha \in \mathbb{F}\),

\begin{equation*} \|\alpha(f,g)\| = \bigl(\|\alpha f\|^{p} + \|\alpha g\|^{p}\bigr)^{1/p} = \bigl(|\alpha|^{p}\|f\|^{p} + |\alpha|^{p}\|g\|^{p}\bigr)^{1/p} = |\alpha| \, \bigl(\|f\|^{p} + \|g\|^{p}\bigr)^{1/p} = |\alpha| \, \|(f,g)\|, \end{equation*}

using homogeneity of the norms on \(V\) and \(W\).

Triangle inequality. Let \(\nu\) denote counting measure on the two-element set \(X = \{1,2\}\), with \(\sigma\)-algebra \(2^{X}\). Every function \(h \colon X \to \mathbb{R}\) is \(2^{X}\)-measurable, and \(|h|^{p}\) is the nonnegative simple function \(|h(1)|^{p} \chi_{\{1\}} + |h(2)|^{p} \chi_{\{2\}}\); because \(\nu(\{1\}) = \nu(\{2\}) = 1\), integration with respect to counting measure is summation (Example 3.6), so

\begin{equation*} \int |h|^{p} \, d\nu = |h(1)|^{p} + |h(2)|^{p} < \infty . \end{equation*}

Hence \(h \in L^{p}(\nu)\) and

\begin{equation*} \|h\|_{p} = \Bigl(\int |h|^{p} \, d\nu\Bigr)^{1/p} = \bigl(|h(1)|^{p} + |h(2)|^{p}\bigr)^{1/p}. \end{equation*}

Now suppose \((f_1, g_1), (f_2, g_2) \in V \times W\). Define \(u, v \colon X \to \mathbb{R}\) by

\begin{equation*} u(1) = \|f_1\|, \quad u(2) = \|g_1\|, \qquad v(1) = \|f_2\|, \quad v(2) = \|g_2\|. \end{equation*}

The triangle inequality in \(V\) gives \(\|f_1 + f_2\| \le \|f_1\| + \|f_2\| = (u+v)(1)\), and similarly \(\|g_1 + g_2\| \le (u+v)(2)\). Because \(t \mapsto t^{p}\) and \(t \mapsto t^{1/p}\) are both increasing on \([0,\infty)\), we have \(\bigl(a_1^{p} + a_2^{p}\bigr)^{1/p} \le \bigl(c_1^{p} + c_2^{p}\bigr)^{1/p}\) whenever \(0 \le a_1 \le c_1\) and \(0 \le a_2 \le c_2\). Hence

\begin{equation*} \begin{aligned} \|(f_1,g_1) + (f_2,g_2)\| &= \bigl(\|f_1+f_2\|^{p} + \|g_1+g_2\|^{p}\bigr)^{1/p} \\ &\le \Bigl(\bigl((u+v)(1)\bigr)^{p} + \bigl((u+v)(2)\bigr)^{p}\Bigr)^{1/p} \\ &= \|u+v\|_{p} \\ &\le \|u\|_{p} + \|v\|_{p} \\ &= \|(f_1,g_1)\| + \|(f_2,g_2)\|, \end{aligned} \end{equation*}

where the second inequality is Minkowski’s inequality (7.14) applied to the measure \(\nu\) and the functions \(u, v \in L^{p}(\nu)\). (For \(p = 1\) this is just the triangle inequality in \(\mathbb{R}\) and needs no appeal to 7.14.) Thus \(\|\cdot\|\) is a norm on \(V \times W\), and \(V \times W\) is a normed vector space.

Completeness. Suppose \((f_1, g_1), (f_2, g_2), \dots\) is a Cauchy sequence in \(V \times W\) with respect to \(\|\cdot\|\). By the first inequality in \((\ast)\) applied to the difference \((f_j - f_k,\, g_j - g_k)\),

\begin{equation*} \|f_j - f_k\| \le \|(f_j, g_j) - (f_k, g_k)\| \quad \text{and} \quad \|g_j - g_k\| \le \|(f_j, g_j) - (f_k, g_k)\| \end{equation*}

for all \(j, k \in \mathbb{Z}^{+}\). Hence given \(\varepsilon > 0\), any \(N\) that works for the Cauchy condition in \(V \times W\) also works for \(f_1, f_2, \dots\) in \(V\) and for \(g_1, g_2, \dots\) in \(W\). So \(f_1, f_2, \dots\) is a Cauchy sequence in \(V\) and \(g_1, g_2, \dots\) is a Cauchy sequence in \(W\). Because \(V\) and \(W\) are Banach spaces, there exist \(f \in V\) and \(g \in W\) with

\begin{equation*} \lim_{k \to \infty} \|f_k - f\| = 0 \quad \text{and} \quad \lim_{k \to \infty} \|g_k - g\| = 0. \end{equation*}

The second inequality in \((\ast)\) now gives

\begin{equation*} \|(f_k, g_k) - (f,g)\| = \|(f_k - f,\, g_k - g)\| \le 2^{1/p} \max\{\|f_k - f\|,\, \|g_k - g\|\} \longrightarrow 0 \end{equation*}

as \(k \to \infty\). Thus \((f_k, g_k) \to (f,g)\) in \(V \times W\). Every Cauchy sequence in \(V \times W\) converges, so \(V \times W\) with this norm is a Banach space.

Remark. The argument above shows more: by \((\ast)\) the norm \(\|\cdot\|\) is equivalent to the norm \(\max\{\|f\|,\|g\|\}\) of 6.84, so the two norms have exactly the same Cauchy sequences and the same convergent sequences (with the same limits). Consequently completeness could also be deduced immediately from 6.84, and, exactly as in 6.84, a sequence \((f_1,g_1), (f_2,g_2), \dots\) converges to \((f,g)\) in the norm above if and only if \(\lim_{k \to \infty} f_k = f\) in \(V\) and \(\lim_{k \to \infty} g_k = g\) in \(W\). Note also that the case \(p = \infty\), interpreted as \(\max\{\|f\|,\|g\|\}\), is precisely 6.84, which is why the exercise restricts to \(1 \le p < \infty\).

Hilbert Spaces

Exercises 8A

Problem (8A.1)

Let \(V\) denote the vector space of bounded continuous functions from \(\mathbb{R}\) to \(\mathbf{F}\). Let \(r_1, r_2, \ldots\) be a list of the rational numbers. For \(f, g \in V\), define

\begin{equation*} \langle f, g \rangle = \sum_{k=1}^{\infty} \frac{f(r_k)\,\overline{g(r_k)}}{2^k}. \end{equation*}

Show that \(\langle \cdot, \cdot \rangle\) is an inner product on \(V\).

Solution

Throughout, for \(h \in V\) write \(\|h\|_\infty = \sup_{x \in \mathbb{R}} |h(x)|\), which is finite because the elements of \(V\) are bounded.

The formula makes sense. Fix \(f, g \in V\). For every \(k \in \mathbb{Z}^+\) we have

\begin{equation*} \left| \frac{f(r_k)\,\overline{g(r_k)}}{2^k} \right| \le \frac{\|f\|_\infty \|g\|_\infty}{2^k}, \end{equation*}

and \(\sum_{k=1}^\infty 2^{-k} = 1\). Hence the series defining \(\langle f, g \rangle\) converges absolutely, so it converges to an element of \(\mathbf{F}\) and \(|\langle f, g \rangle| \le \|f\|_\infty \|g\|_\infty\). Thus \(\langle \cdot, \cdot \rangle\) is a well-defined function from \(V \times V\) to \(\mathbf{F}\). Because all the series appearing below converge absolutely, the manipulations that follow (splitting a sum into two sums, pulling a scalar out of a sum) are legitimate.

We now verify the four properties in the definition 8.1 of an inner product.

Positivity. For \(f \in V\),

\begin{equation*} \langle f, f \rangle = \sum_{k=1}^{\infty} \frac{f(r_k)\,\overline{f(r_k)}}{2^k} = \sum_{k=1}^{\infty} \frac{|f(r_k)|^2}{2^k}, \end{equation*}

which is a convergent series of nonnegative real numbers and therefore lies in \([0, \infty)\).

Definiteness. If \(f = 0\) then clearly \(\langle f, f \rangle = 0\). Conversely, suppose \(f \in V\) and \(\langle f, f \rangle = 0\). The displayed formula above exhibits \(\langle f, f \rangle\) as a sum of nonnegative terms; a sum of nonnegative numbers is \(0\) only if every term is \(0\). Hence \(|f(r_k)|^2 / 2^k = 0\), and thus \(f(r_k) = 0\), for every \(k \in \mathbb{Z}^+\). In other words \(f\) vanishes at every rational number. Now let \(x \in \mathbb{R}\). Because \(\mathbb{Q}\) is dense in \(\mathbb{R}\), there is a sequence \((q_n)\) of rational numbers with \(q_n \to x\). Continuity of \(f\) gives

\begin{equation*} f(x) = \lim_{n \to \infty} f(q_n) = \lim_{n \to \infty} 0 = 0. \end{equation*}

Hence \(f = 0\), as required. (This is the only place where continuity is used, and it is genuinely needed: without it, the function equal to \(0\) on \(\mathbb{Q}\) and \(1\) elsewhere would be a nonzero element with \(\langle f, f \rangle = 0\).)

Linearity in the first slot. Suppose \(f, g, h \in V\) and \(\alpha \in \mathbf{F}\). Then

\begin{equation*} \begin{aligned} \langle f + g, h \rangle &= \sum_{k=1}^{\infty} \frac{\bigl(f(r_k) + g(r_k)\bigr)\overline{h(r_k)}}{2^k} \\ &= \sum_{k=1}^{\infty} \frac{f(r_k)\overline{h(r_k)}}{2^k} + \sum_{k=1}^{\infty} \frac{g(r_k)\overline{h(r_k)}}{2^k} \\ &= \langle f, h \rangle + \langle g, h \rangle, \end{aligned} \end{equation*}

where the middle equality splits the sum into two series, each of which converges absolutely by the estimate at the start of this solution. Similarly,

\begin{equation*} \langle \alpha f, h \rangle = \sum_{k=1}^{\infty} \frac{\alpha f(r_k)\overline{h(r_k)}}{2^k} = \alpha \sum_{k=1}^{\infty} \frac{f(r_k)\overline{h(r_k)}}{2^k} = \alpha \langle f, h \rangle. \end{equation*}

Conjugate symmetry. Suppose \(f, g \in V\). Complex conjugation is continuous and additive, so it commutes with the limit of the partial sums of an absolutely convergent series. Hence

\begin{equation*} \overline{\langle g, f \rangle} = \overline{\sum_{k=1}^{\infty} \frac{g(r_k)\overline{f(r_k)}}{2^k}} = \sum_{k=1}^{\infty} \overline{\left(\frac{g(r_k)\overline{f(r_k)}}{2^k}\right)} = \sum_{k=1}^{\infty} \frac{f(r_k)\overline{g(r_k)}}{2^k} = \langle f, g \rangle, \end{equation*}

using that \(2^k\) is real.

All four properties hold, so \(\langle \cdot, \cdot \rangle\) is an inner product on \(V\).

Problem (8A.2)

Prove that if \(\mu\) is a measure and \(f, g \in L^2(\mu)\), then

\begin{equation*} \|f\|^2 \|g\|^2 - |\langle f, g \rangle|^2 = \frac{1}{2} \int \int |f(x) g(y) - g(x) f(y)|^2 \, d\mu(y) \, d\mu(x). \end{equation*}

Solution

Here \(\|\cdot\|\) and \(\langle \cdot, \cdot \rangle\) are the norm and the standard inner product on \(L^2(\mu)\), so \(\|f\|^2 = \int |f|^2 \, d\mu\) and \(\langle f, g \rangle = \int f \overline{g} \, d\mu\) (see 8.2 and 8.5). Write

\begin{equation*} A = \int |f|^2 \, d\mu = \|f\|^2, \qquad B = \int |g|^2 \, d\mu = \|g\|^2, \qquad C = \int f \overline{g} \, d\mu = \langle f, g \rangle . \end{equation*}

Note \(A, B < \infty\) because \(f, g \in L^2(\mu)\), and \(f \overline{g} \in L^1(\mu)\) with \(|C| \le \sqrt{A}\sqrt{B}\) by the Cauchy–Schwarz inequality for \(L^2(\mu)\) (8.14). Also \(\int g \overline{f} \, d\mu = \overline{\langle f, g \rangle} = \overline{C}\) by conjugate symmetry.

We evaluate the right side as an iterated integral, exactly as it is written: first integrate in \(y\) with \(x\) fixed, then integrate the result in \(x\). This requires only linearity of the integral, so no appeal to Tonelli’s theorem (and hence no \(\sigma\)-finiteness hypothesis) is needed.

The inner integral. Let \((X, \mathcal{S}, \mu)\) be the underlying measure space and fix \(x \in X\). Because \(f\) and \(g\) take values in \(\mathbf{F}\), the numbers \(f(x)\) and \(g(x)\) are finite scalars, so the computation below is valid for every \(x \in X\), not merely almost every \(x\); this matters, since it lets us conclude that the function of \(x\) produced below is genuinely \(\mathcal{S}\)-measurable rather than only almost equal to a measurable function (we are not assuming \(\mu\) is a complete measure). Expanding \(|w|^2 = w \overline{w}\) with \(w = f(x)g(y) - g(x)f(y)\) gives, for each \(y\),

\begin{equation*} \begin{aligned} |f(x)g(y) - g(x)f(y)|^2 &= |f(x)|^2 |g(y)|^2 + |g(x)|^2 |f(y)|^2 \\ &\qquad - f(x)\overline{g(x)}\, g(y)\overline{f(y)} - \overline{f(x)}\,g(x)\, f(y)\overline{g(y)} , \end{aligned} \end{equation*}

where the two cross terms have been rearranged so that the \(x\)-dependent scalars are separated from the \(y\)-dependent functions. As a function of \(y\), the right side is a linear combination, with the finite scalar coefficients \(|f(x)|^2\), \(|g(x)|^2\), \(-f(x)\overline{g(x)}\), \(-\overline{f(x)}g(x)\), of the four functions \(|g|^2\), \(|f|^2\), \(g\overline{f}\), \(f\overline{g}\), each of which is in \(L^1(\mu)\) as noted above. Hence the integrand is \(\mu\)-integrable in \(y\), and by linearity of the integral

\begin{equation*} \int |f(x)g(y) - g(x)f(y)|^2 \, d\mu(y) = |f(x)|^2 B + |g(x)|^2 A - f(x)\overline{g(x)}\,\overline{C} - \overline{f(x)}g(x)\, C . \end{equation*}

The outer integral. The right side of the last display is, as a function of \(x\), a linear combination with finite scalar coefficients \(B\), \(A\), \(-\overline{C}\), \(-C\) of the functions \(|f|^2\), \(|g|^2\), \(f\overline{g}\), \(\overline{f}g\), all in \(L^1(\mu)\). So it is integrable in \(x\) and, again by linearity,

\begin{equation*} \begin{aligned} \int \int |f(x)g(y) - g(x)f(y)|^2 \, d\mu(y) \, d\mu(x) &= A B + B A - \overline{C} \int f \overline{g} \, d\mu - C \int \overline{f} g \, d\mu \\ &= 2AB - \overline{C}\, C - C\, \overline{C} \\ &= 2AB - 2|C|^2 . \end{aligned} \end{equation*}

Multiplying by \(\tfrac{1}{2}\) gives

\begin{equation*} \frac{1}{2}\int \int |f(x)g(y) - g(x)f(y)|^2 \, d\mu(y) \, d\mu(x) = AB - |C|^2 = \|f\|^2 \|g\|^2 - |\langle f, g \rangle|^2, \end{equation*}

as desired.

Remark. The right side of the identity is \(\tfrac12\) times an iterated integral of a nonnegative function, hence is nonnegative. So the identity yields \(\|f\|^2\|g\|^2 - |\langle f, g\rangle|^2 \ge 0\), which is the Cauchy–Schwarz inequality \(|\langle f, g\rangle| \le \|f\|\,\|g\|\) on \(L^2(\mu)\) (8.11) — proved here without the orthogonal decomposition 8.10. The identity also shows that equality holds in Cauchy–Schwarz for these particular \(f, g\) if and only if the iterated integral on the right vanishes.

Problem (8A.3)

Suppose \(f\) and \(g\) are elements of an inner product space and

\begin{equation*} \|f + g\|^2 = \|f\|^2 + \|g\|^2 . \end{equation*}

(a) Prove that if \(\mathbf{F} = \mathbb{R}\), then \(f\) and \(g\) are orthogonal.

(b) Give an example to show that if \(\mathbf{F} = \mathbb{C}\), then \(f\) and \(g\) can satisfy the equation above without being orthogonal.

Solution

First we expand \(\|f+g\|^2\) in any inner product space. Using the definition 8.4 of the norm and then linearity in the first slot together with 8.3(b),

\begin{equation*} \begin{aligned} \|f+g\|^2 &= \langle f + g, f + g \rangle \\ &= \langle f, f \rangle + \langle g, f \rangle + \langle f, g \rangle + \langle g, g \rangle \\ &= \|f\|^2 + \|g\|^2 + \langle f, g \rangle + \overline{\langle f, g \rangle} \\ &= \|f\|^2 + \|g\|^2 + 2 \operatorname{Re} \langle f, g \rangle, \end{aligned} \end{equation*}

where the third equality uses conjugate symmetry, \(\langle g, f \rangle = \overline{\langle f, g \rangle}\). (This is the computation appearing in the proof of the triangle inequality 8.15.) Hence the hypothesis \(\|f+g\|^2 = \|f\|^2 + \|g\|^2\) is equivalent to

\begin{equation*} \operatorname{Re}\langle f, g \rangle = 0 . \end{equation*}

(a) Suppose \(\mathbf{F} = \mathbb{R}\). Then \(\langle f, g \rangle\) is a real number, so \(\langle f, g \rangle = \operatorname{Re}\langle f, g \rangle = 0\). By the definition 8.7 of orthogonality, \(f\) and \(g\) are orthogonal. Thus for real inner product spaces the converse of the Pythagorean Theorem 8.9 holds.

(b) Suppose \(\mathbf{F} = \mathbb{C}\). Take the inner product space \(V = \mathbb{C}\) with its standard inner product \(\langle w, z \rangle = w \overline{z}\) (the case \(n = 1\) of 8.2), and let

\begin{equation*} f = 1, \qquad g = i . \end{equation*}

Then \(\|f\|^2 = 1\) and \(\|g\|^2 = |i|^2 = 1\), while

\begin{equation*} \|f + g\|^2 = |1 + i|^2 = 1^2 + 1^2 = 2 = \|f\|^2 + \|g\|^2 , \end{equation*}

so the displayed equation in the exercise holds. But

\begin{equation*} \langle f, g \rangle = 1 \cdot \overline{i} = -i \ne 0, \end{equation*}

so \(f\) and \(g\) are not orthogonal. (Consistently with the computation above, \(\operatorname{Re}(-i) = 0\).)

The same example transplants to any complex inner product space containing a nonzero element \(u\): take \(f = u\) and \(g = i u\). Then \(\langle f, g \rangle = \langle u, iu \rangle = \overline{i}\,\|u\|^2 = -i \|u\|^2 \ne 0\) by 8.3(c), while \(\|f + g\|^2 = \|(1+i)u\|^2 = |1+i|^2 \|u\|^2 = 2\|u\|^2 = \|f\|^2 + \|g\|^2\) by 8.6.

Problem (8A.4)

Find \(a, b \in \mathbb{R}^3\) such that \(a\) is a scalar multiple of \((1, 6, 3)\), \(b\) is orthogonal to \((1, 6, 3)\), and \((5, 4, -2) = a + b\).

Solution

Write \(u = (5, 4, -2)\) and \(v = (1, 6, 3)\), and use the standard inner product on \(\mathbb{R}^3\) from 8.2, namely \(\langle (a_1,a_2,a_3),(b_1,b_2,b_3)\rangle = a_1 b_1 + a_2 b_2 + a_3 b_3\).

The orthogonal decomposition 8.10 (with \(f = u\) and \(g = v \ne 0\)) tells us exactly what to take: \(a = \frac{\langle u, v\rangle}{\|v\|^2} v\) and \(b = u - a\). Compute

\begin{equation*} \langle u, v \rangle = 5 \cdot 1 + 4 \cdot 6 + (-2)\cdot 3 = 5 + 24 - 6 = 23, \end{equation*}

\begin{equation*} \|v\|^2 = 1^2 + 6^2 + 3^2 = 1 + 36 + 9 = 46 . \end{equation*}

So the required scalar is \(\frac{23}{46} = \frac{1}{2}\), and we set

\begin{equation*} a = \tfrac{1}{2}(1, 6, 3) = \left(\tfrac{1}{2},\, 3,\, \tfrac{3}{2}\right), \qquad b = (5, 4, -2) - \left(\tfrac{1}{2},\, 3,\, \tfrac{3}{2}\right) = \left(\tfrac{9}{2},\, 1,\, -\tfrac{7}{2}\right). \end{equation*}

We verify the three requirements directly.

First, \(a = \tfrac{1}{2}(1,6,3)\) is by construction a scalar multiple of \((1,6,3)\).

Second, \(b\) is orthogonal to \((1,6,3)\), because

\begin{equation*} \left\langle \left(\tfrac{9}{2}, 1, -\tfrac{7}{2}\right), (1, 6, 3) \right\rangle = \tfrac{9}{2} + 6 - \tfrac{21}{2} = \tfrac{9 - 21}{2} + 6 = -6 + 6 = 0 . \end{equation*}

Third, \(a + b = \left(\tfrac{1}{2} + \tfrac{9}{2},\, 3 + 1,\, \tfrac{3}{2} - \tfrac{7}{2}\right) = (5, 4, -2)\), as required.

Uniqueness. This is the only such pair. Indeed, if \(a = c(1,6,3)\) with \(c \in \mathbb{R}\) and \(b = (5,4,-2) - a\) is orthogonal to \((1,6,3)\), then taking the inner product of \((5,4,-2) = c(1,6,3) + b\) with \((1,6,3)\) gives \(23 = 46c + 0\), so \(c = \tfrac{1}{2}\), which forces \(a\) and hence \(b\) to be as above.

Problem (8A.5)

Prove that

\begin{equation*} 16 \le (a + b + c + d)\left(\frac{1}{a} + \frac{1}{b} + \frac{1}{c} + \frac{1}{d}\right) \end{equation*}

for all positive numbers \(a, b, c, d\), with equality if and only if \(a = b = c = d\).

Solution

Let \(a, b, c, d\) be positive numbers. Work in the inner product space \(\mathbb{R}^4\) with its standard inner product (8.2), and set

\begin{equation*} u = \left(\sqrt{a},\ \sqrt{b},\ \sqrt{c},\ \sqrt{d}\right), \qquad v = \left(\frac{1}{\sqrt{a}},\ \frac{1}{\sqrt{b}},\ \frac{1}{\sqrt{c}},\ \frac{1}{\sqrt{d}}\right). \end{equation*}

These are well defined because \(a, b, c, d > 0\), and both are nonzero vectors. Their inner product is

\begin{equation*} \langle u, v \rangle = \sqrt{a}\cdot\frac{1}{\sqrt{a}} + \sqrt{b}\cdot\frac{1}{\sqrt{b}} + \sqrt{c}\cdot\frac{1}{\sqrt{c}} + \sqrt{d}\cdot\frac{1}{\sqrt{d}} = 4, \end{equation*}

while

\begin{equation*} \|u\|^2 = a + b + c + d, \qquad \|v\|^2 = \frac{1}{a} + \frac{1}{b} + \frac{1}{c} + \frac{1}{d}. \end{equation*}

The Cauchy–Schwarz inequality 8.11 gives \(|\langle u, v \rangle| \le \|u\|\,\|v\|\); squaring (both sides are nonnegative) yields

\begin{equation*} 16 = |\langle u, v\rangle|^2 \le \|u\|^2 \|v\|^2 = (a + b + c + d)\left(\frac{1}{a} + \frac{1}{b} + \frac{1}{c} + \frac{1}{d}\right), \end{equation*}

which is the desired inequality.

Equality. By 8.11, equality holds if and only if one of \(u, v\) is a scalar multiple of the other. Since \(u \ne 0\) and \(v \ne 0\), either form of that statement implies \(v = \lambda u\) for some \(\lambda \in \mathbb{R}\) with \(\lambda \ne 0\). (If instead \(u = \mu v\), then \(\mu \ne 0\) because \(u \ne 0\), and we may take \(\lambda = 1/\mu\).)

Suppose \(v = \lambda u\). Comparing first coordinates, \(\frac{1}{\sqrt{a}} = \lambda \sqrt{a}\), so \(\lambda = \frac{1}{a}\). The same comparison in the other three coordinates gives \(\lambda = \frac{1}{b} = \frac{1}{c} = \frac{1}{d}\) as well. Hence \(\frac1a = \frac1b = \frac1c = \frac1d\), and since \(a,b,c,d\) are positive this means \(a = b = c = d\).

Conversely, if \(a = b = c = d\), then

\begin{equation*} (a+b+c+d)\left(\frac{1}{a} + \frac{1}{b} + \frac{1}{c} + \frac{1}{d}\right) = (4a)\cdot\frac{4}{a} = 16, \end{equation*}

so equality holds. Therefore equality holds if and only if \(a = b = c = d\).

Problem (8A.6)

Prove that the square of the average of each finite list of real numbers containing at least two distinct real numbers is less than the average of the squares of the numbers in that list.

Solution

Let \(x_1, \ldots, x_n\) be a finite list of real numbers containing at least two distinct real numbers; in particular \(n \ge 2\). We must show

\begin{equation*} \left(\frac{x_1 + \cdots + x_n}{n}\right)^{2} < \frac{x_1^2 + \cdots + x_n^2}{n} . \end{equation*}

Work in the inner product space \(\mathbb{R}^n\) with its standard inner product (8.2), and set

\begin{equation*} u = (x_1, \ldots, x_n), \qquad v = (1, 1, \ldots, 1). \end{equation*}

Then

\begin{equation*} \langle u, v \rangle = x_1 + \cdots + x_n, \qquad \|u\|^2 = x_1^2 + \cdots + x_n^2, \qquad \|v\|^2 = n . \end{equation*}

Squaring the Cauchy–Schwarz inequality 8.11, \(|\langle u, v\rangle| \le \|u\|\,\|v\|\), gives

\begin{equation*} (x_1 + \cdots + x_n)^2 \le n\,(x_1^2 + \cdots + x_n^2). \end{equation*}

Dividing both sides by \(n^2 > 0\) yields

\begin{equation*} \left(\frac{x_1 + \cdots + x_n}{n}\right)^{2} \le \frac{x_1^2 + \cdots + x_n^2}{n} . \end{equation*}

It remains to rule out equality. Equality in the display above holds if and only if it holds in the Cauchy–Schwarz inequality, which by 8.11 happens if and only if one of \(u, v\) is a scalar multiple of the other. Since \(v \ne 0\), that condition forces \(u\) to be a scalar multiple of \(v\): if \(u = \lambda v\) we are done, and if instead \(v = \mu u\) then \(\mu \ne 0\) (as \(v \ne 0\)) and \(u = \frac{1}{\mu} v\). But \(u = \lambda v\) means

\begin{equation*} x_1 = x_2 = \cdots = x_n = \lambda, \end{equation*}

that is, all the numbers in the list are equal. This contradicts the hypothesis that the list contains at least two distinct real numbers. Hence equality cannot hold, and we conclude

\begin{equation*} \left(\frac{x_1 + \cdots + x_n}{n}\right)^{2} < \frac{x_1^2 + \cdots + x_n^2}{n}, \end{equation*}

as desired.

Remark. Equivalently, this says that the variance \(\frac1n\sum_{k=1}^n x_k^2 - \bigl(\frac1n \sum_{k=1}^n x_k\bigr)^2\) of a list of real numbers is strictly positive unless all entries of the list are equal.

Problem (8A.7)

Suppose \(f\) and \(g\) are elements of an inner product space and \(\|f\| \le 1\) and \(\|g\| \le 1\). Prove that

\begin{equation*} \sqrt{1 - \|f\|^2}\,\sqrt{1 - \|g\|^2} \le 1 - |\langle f, g \rangle| . \end{equation*}

Solution

Write

\begin{equation*} s = \|f\|, \qquad t = \|g\| . \end{equation*}

Then \(s, t \in [0, 1]\): indeed \(\|h\| = \sqrt{\langle h, h\rangle} \ge 0\) for every \(h\), by the definition 8.4 of the norm together with the positivity property in the definition 8.1 of an inner product (this is also recorded in the discussion following 8.5), and by hypothesis \(s \le 1\) and \(t \le 1\). In particular \(1 - s^2 \ge 0\) and \(1 - t^2 \ge 0\), so the square roots in the statement are square roots of nonnegative real numbers and everything below is well defined.

Step 1: reduce to an inequality about \(s\) and \(t\). By the Cauchy–Schwarz inequality 8.11,

\begin{equation*} |\langle f, g \rangle| \le \|f\|\,\|g\| = st . \end{equation*}

Hence

\begin{equation*} 1 - |\langle f, g \rangle| \ge 1 - st, \end{equation*}

so it suffices to prove the numerical inequality

\begin{equation*} \sqrt{1 - s^2}\,\sqrt{1 - t^2} \le 1 - st \end{equation*}

for all \(s, t \in [0,1]\). Call this inequality \((\ast)\).

Step 2: prove \((\ast)\). Both sides of \((\ast)\) are nonnegative: the left side is a product of two nonnegative square roots, and the right side satisfies \(1 - st \ge 1 - 1 \cdot 1 = 0\) because \(0 \le s \le 1\) and \(0 \le t \le 1\). For nonnegative real numbers, the inequality \(\alpha \le \beta\) holds if and only if \(\alpha^2 \le \beta^2\). So \((\ast)\) is equivalent to

\begin{equation*} (1 - s^2)(1 - t^2) \le (1 - st)^2 . \end{equation*}

Expanding both sides,

\begin{equation*} 1 - s^2 - t^2 + s^2 t^2 \le 1 - 2st + s^2 t^2 , \end{equation*}

which after cancelling \(1\) and \(s^2 t^2\) from both sides is equivalent to

\begin{equation*} -s^2 - t^2 \le -2st, \qquad \text{i.e.} \qquad 2st \le s^2 + t^2, \qquad \text{i.e.} \qquad 0 \le (s - t)^2 . \end{equation*}

The last inequality is true for all real \(s, t\). Reversing the chain of equivalences proves \((\ast)\).

Conclusion. Combining the two steps,

\begin{equation*} \sqrt{1 - \|f\|^2}\,\sqrt{1 - \|g\|^2} = \sqrt{1 - s^2}\,\sqrt{1 - t^2} \le 1 - st \le 1 - |\langle f, g \rangle| , \end{equation*}

as desired.

Remark on equality. Tracing the two steps, equality holds if and only if \(s = t\) (equality in \(0 \le (s-t)^2\)) and \(|\langle f, g\rangle| = st\) (equality in Cauchy–Schwarz). Thus equality occurs exactly when \(\|f\| = \|g\|\) and one of \(f, g\) is a scalar multiple of the other; equivalently, when \(g = \alpha f\) for some \(\alpha \in \mathbf{F}\) with \(|\alpha| = 1\), or \(f = g = 0\).

Problem (8A.8)

Suppose \(a\) and \(b\) are nonzero elements of \(\mathbb{R}^2\). Prove that

\begin{equation*} \langle a, b\rangle = \|a\|\,\|b\|\cos\theta, \end{equation*}

where \(\theta\) is the angle between \(a\) and \(b\) (thinking of \(a\) as the vector whose initial point is the origin and whose end point is \(a\), and similarly for \(b\)).

Hint: Draw the triangle formed by \(a\), \(b\), and \(a - b\); then use the law of cosines.

Solution

Throughout, \(\mathbb{R}^2\) carries the standard inner product, so that \(\langle x,y\rangle = \langle y,x\rangle\) for all \(x,y \in \mathbb{R}^2\) and \(\|x\| = \sqrt{\langle x,x\rangle}\) is the usual Euclidean length (see 8.5). The angle \(\theta\) between \(a\) and \(b\) is by definition the angle at the origin between the ray through \(a\) and the ray through \(b\); thus \(\theta \in [0,\pi]\).

First we compute \(\|a-b\|^2\) algebraically. Expanding by additivity in each slot and using symmetry of the real inner product,

\begin{equation*} \begin{aligned} \|a-b\|^2 &= \langle a-b,\, a-b\rangle\\ &= \langle a,a\rangle - \langle a,b\rangle - \langle b,a\rangle + \langle b,b\rangle\\ &= \|a\|^2 + \|b\|^2 - 2\langle a,b\rangle. \end{aligned} \end{equation*}

Now we compute \(\|a-b\|^2\) geometrically. Place the initial points of \(a\) and \(b\) at the origin \(O\), and let \(A\) and \(B\) denote their end points. The segment from \(B\) to \(A\) is the arrow representing \(a - b\), so the triangle \(OAB\) has side lengths

\begin{equation*} |OA| = \|a\|,\qquad |OB| = \|b\|,\qquad |AB| = \|a-b\|, \end{equation*}

and the angle of the triangle at the vertex \(O\) is exactly the angle \(\theta\) between \(a\) and \(b\). The law of cosines applied to this triangle gives

\begin{equation*} \|a-b\|^2 = \|a\|^2 + \|b\|^2 - 2\|a\|\,\|b\|\cos\theta . \end{equation*}

Comparing the two displayed expressions for \(\|a-b\|^2\) and cancelling \(\|a\|^2 + \|b\|^2\) from both sides yields

\begin{equation*} -2\langle a,b\rangle = -2\|a\|\,\|b\|\cos\theta, \end{equation*}

and dividing by \(-2\) gives \(\langle a,b\rangle = \|a\|\,\|b\|\cos\theta\), as desired.

It remains to treat the degenerate cases, in which \(O\), \(A\), \(B\) are collinear and the triangle above collapses (so that appealing to the law of cosines is not legitimate). This happens exactly when \(b\) is a real scalar multiple of \(a\); write \(b = ta\) with \(t \in \mathbb{R}\), \(t \neq 0\), since \(a\) and \(b\) are nonzero.

If \(t > 0\), then \(a\) and \(b\) point along the same ray from the origin, so \(\theta = 0\) and \(\cos\theta = 1\). Here \(\langle a,b\rangle = t\|a\|^2 = \|a\|\,(t\|a\|) = \|a\|\,\|b\|\), using \(\|b\| = |t|\,\|a\| = t\|a\|\) by 8.6. Thus \(\langle a,b\rangle = \|a\|\,\|b\|\cos\theta\).

If \(t < 0\), then \(a\) and \(b\) point along opposite rays, so \(\theta = \pi\) and \(\cos\theta = -1\). Here \(\langle a,b\rangle = t\|a\|^2 = -|t|\,\|a\|^2 = -\|a\|\,\|b\|\), again by 8.6. Thus once more \(\langle a,b\rangle = \|a\|\,\|b\|\cos\theta\).

Hence the identity holds for all nonzero \(a, b \in \mathbb{R}^2\).

Problem (8A.9)

The angle between two vectors (thought of as arrows with initial point at the origin) in \(\mathbb{R}^2\) or \(\mathbb{R}^3\) can be defined geometrically. However, geometry is not as clear in \(\mathbb{R}^n\) for \(n > 3\). Thus the angle between two nonzero vectors \(a, b \in \mathbb{R}^n\) is defined to be

\begin{equation*} \arccos \frac{\langle a,b\rangle}{\|a\|\,\|b\|}, \end{equation*}

where the motivation for this definition comes from the previous exercise. Explain why the Cauchy–Schwarz inequality is needed to show that this definition makes sense.

Solution

The function \(\cos\) is a bijection from \([0,\pi]\) onto \([-1,1]\), and \(\arccos\) is by definition the inverse of that bijection. Hence the domain of \(\arccos\) is precisely the interval \([-1,1]\): the expression \(\arccos t\) is meaningful (as a real number) only when \(t \in [-1,1]\), because no real number has cosine outside \([-1,1]\).

So for the proposed definition to make sense, we must know that

\begin{equation*} \frac{\langle a,b\rangle}{\|a\|\,\|b\|} \in [-1,1] \end{equation*}

for all nonzero \(a,b \in \mathbb{R}^n\). Since \(a\) and \(b\) are nonzero we have \(\|a\|\,\|b\| > 0\), so the quotient is a well-defined real number, and the requirement above is equivalent to

\begin{equation*} |\langle a,b\rangle| \le \|a\|\,\|b\| . \end{equation*}

That is exactly the Cauchy–Schwarz inequality (8.11), applied in the inner product space \(\mathbb{R}^n\). Without it there would be no reason at all to expect the quotient to lie in \([-1,1]\), and the formula would define nothing.

Thus the Cauchy–Schwarz inequality is precisely the statement that guarantees the quotient lies in the domain of \(\arccos\); granted it, the definition assigns to each pair of nonzero vectors \(a,b \in \mathbb{R}^n\) a unique angle in \([0,\pi]\), which is the expected range of values for the angle between two vectors.

Two further remarks show that the definition behaves as geometry demands, and both again rest on 8.11. First, by 8A.8 the definition agrees with the geometric notion of angle when \(n = 2\) (and, by the same law-of-cosines argument carried out in the plane containing \(a\), \(b\), and the origin, when \(n = 3\)); so it is a genuine extension of the geometric definition rather than a competing one. Second, the equality condition in 8.11 says that \(|\langle a,b\rangle| = \|a\|\,\|b\|\) precisely when one of \(a\), \(b\) is a scalar multiple of the other; hence the angle equals \(0\) or \(\pi\) (the two extreme values, attained only at the endpoints of \([-1,1]\)) exactly when \(a\) and \(b\) lie on the same or on opposite rays through the origin, as it should.

Problem (8A.10)

(a) Suppose \(f\) and \(g\) are elements of a real inner product space. Prove that \(f\) and \(g\) have the same norm if and only if \(f + g\) is orthogonal to \(f - g\).

(b) Use (a) to show that the diagonals of a parallelogram are perpendicular to each other if and only if the parallelogram is a rhombus.

Solution

(a) Because the inner product space is real, the conjugate symmetry property of an inner product (8.1) reads \(\langle u,v\rangle = \langle v,u\rangle\) for all \(u,v\) in the space. Using additivity in the first slot (part of 8.1) together with 8.3(b) and 8.3(c) for the second slot (so that \(\langle u,-v\rangle = -\langle u,v\rangle\)), we compute

\begin{equation*} \begin{aligned} \langle f+g,\, f-g\rangle &= \langle f,f\rangle - \langle f,g\rangle + \langle g,f\rangle - \langle g,g\rangle\\ &= \|f\|^2 - \langle f,g\rangle + \langle f,g\rangle - \|g\|^2\\ &= \|f\|^2 - \|g\|^2 . \end{aligned} \end{equation*}

By 8.7, \(f+g\) is orthogonal to \(f-g\) if and only if \(\langle f+g, f-g\rangle = 0\), which by the computation above holds if and only if \(\|f\|^2 = \|g\|^2\). Since norms are nonnegative and \(t \mapsto t^2\) is injective on \([0,\infty)\), this last equation holds if and only if \(\|f\| = \|g\|\). Hence \(f\) and \(g\) have the same norm if and only if \(f+g\) is orthogonal to \(f-g\).

Note that the hypothesis that the scalar field is \(\mathbb{R}\) is essential: in a complex inner product space the middle line above becomes \(\|f\|^2 - \|g\|^2 + \big(\overline{\langle f,g\rangle} - \langle f,g\rangle\big) = \|f\|^2 - \|g\|^2 - 2i\operatorname{Im}\langle f,g\rangle\), and the imaginary part need not vanish. For instance in \(\mathbb{C}\) (a one-dimensional complex inner product space, with \(\langle z,w\rangle = z\overline{w}\)) take \(f = 1\) and \(g = i\): then \(\|f\| = \|g\| = 1\), yet \(\langle f+g, f-g\rangle = (1+i)\overline{(1-i)} = (1+i)^2 = 2i \neq 0\).

(b) Work in \(\mathbb{R}^2\) with its standard inner product. As explained in the discussion following 8.8 (which rests on 8A.8), two vectors in \(\mathbb{R}^2\) are orthogonal precisely when they are perpendicular in the usual sense of plane geometry, so the two words may be used interchangeably here. Perpendicularity of segments and equality of side lengths are unchanged by translation, so we may position the parallelogram with one vertex at the origin. Its four vertices are then

\begin{equation*} 0,\qquad f,\qquad f+g,\qquad g, \end{equation*}

where \(f\) and \(g\) are the two nonzero, non-parallel vectors spanning the parallelogram; the four sides are the segments \(0f\), \(f\,(f+g)\), \((f+g)\,g\), and \(g\,0\), which as arrows are \(f\), \(g\), \(-f\), \(-g\) respectively.

The two diagonals join \(0\) to \(f+g\) and join \(f\) to \(g\). As arrows these are \(f+g\) and \(g-f\), and two segments are perpendicular exactly when the arrows along them are orthogonal. By 8.3(c) with \(\alpha = -1\) (a real scalar, so equal to its own conjugate),

\begin{equation*} \langle f+g,\, g-f\rangle = \langle f+g,\, -(f-g)\rangle = -\langle f+g,\, f-g\rangle, \end{equation*}

so \(g-f\) is orthogonal to \(f+g\) if and only if \(f-g\) is orthogonal to \(f+g\). Hence the diagonals are perpendicular if and only if \(\langle f+g,\, f-g\rangle = 0\), which by part (a) holds if and only if \(\|f\| = \|g\|\).

Finally, the side lengths of the parallelogram are \(\|f\|\), \(\|g\|\), \(\|-f\| = \|f\|\), and \(\|-g\| = \|g\|\) (by 8.6, homogeneity of the norm, with \(\alpha = -1\)). Hence all four sides have equal length if and only if \(\|f\| = \|g\|\), i.e. if and only if the parallelogram is a rhombus. Combining this with the chain of equivalences above, the diagonals of the parallelogram are perpendicular to each other if and only if the parallelogram is a rhombus.

Problem (8A.11)

Suppose \(f\) and \(g\) are elements of an inner product space. Prove that \(\|f\| = \|g\|\) if and only if \(\|sf + tg\| = \|tf + sg\|\) for all \(s, t \in \mathbb{R}\).

Solution

The inner product space may be real or complex; the scalars \(s,t\) in the statement are real, which is what makes the computation below work in both cases.

Fix \(s, t \in \mathbb{R}\). Expanding \(\|sf+tg\|^2 = \langle sf+tg,\, sf+tg\rangle\) using additivity in each slot, homogeneity in the first slot, and 8.3(c) for the second slot (where the conjugates of the real numbers \(s,t\) are \(s,t\) themselves), we get

\begin{equation*} \begin{aligned} \|sf+tg\|^2 &= s^2\langle f,f\rangle + st\langle f,g\rangle + ts\langle g,f\rangle + t^2\langle g,g\rangle\\ &= s^2\|f\|^2 + t^2\|g\|^2 + st\big(\langle f,g\rangle + \overline{\langle f,g\rangle}\big)\\ &= s^2\|f\|^2 + t^2\|g\|^2 + 2st\operatorname{Re}\langle f,g\rangle, \end{aligned} \end{equation*}

where the conjugate symmetry property of an inner product gave \(\langle g,f\rangle = \overline{\langle f,g\rangle}\). Interchanging the roles of \(s\) and \(t\) in this identity (the term \(2st\operatorname{Re}\langle f,g\rangle\) is unchanged, since \(st = ts\)) gives

\begin{equation*} \|tf+sg\|^2 = t^2\|f\|^2 + s^2\|g\|^2 + 2st\operatorname{Re}\langle f,g\rangle . \end{equation*}

Subtracting the second display from the first, the cross terms cancel and we obtain the key identity

\begin{equation*} \|sf+tg\|^2 - \|tf+sg\|^2 = (s^2 - t^2)\big(\|f\|^2 - \|g\|^2\big) \end{equation*}

for all \(s,t \in \mathbb{R}\).

Now suppose \(\|f\| = \|g\|\). Then the right side of the key identity is \(0\) for every choice of \(s,t \in \mathbb{R}\), so \(\|sf+tg\|^2 = \|tf+sg\|^2\); taking nonnegative square roots gives \(\|sf+tg\| = \|tf+sg\|\) for all \(s,t \in \mathbb{R}\).

Conversely, suppose \(\|sf+tg\| = \|tf+sg\|\) for all \(s,t \in \mathbb{R}\). Take \(s = 1\) and \(t = 0\). Then the hypothesis reads \(\|f\| = \|g\|\), as desired. (Alternatively, apply the key identity with \(s = 1\), \(t = 0\): its left side is \(0\) by hypothesis and its right side is \(\|f\|^2 - \|g\|^2\), so \(\|f\|^2 = \|g\|^2\) and hence \(\|f\| = \|g\|\) since norms are nonnegative.)

Thus \(\|f\| = \|g\|\) if and only if \(\|sf+tg\| = \|tf+sg\|\) for all \(s,t \in \mathbb{R}\).

Problem (8A.12)

Suppose \(f\) and \(g\) are elements of an inner product space and \(\|f\| = \|g\| = 1\) and \(\langle f,g\rangle = 1\). Prove that \(f = g\).

Solution

Because \(\langle f,g\rangle = 1\) and \(1\) is a real number, conjugate symmetry gives

\begin{equation*} \langle g,f\rangle = \overline{\langle f,g\rangle} = \overline{1} = 1 . \end{equation*}

Now expand \(\|f-g\|^2\), using additivity in each slot and 8.3(c) with \(\alpha = -1\):

\begin{equation*} \begin{aligned} \|f-g\|^2 &= \langle f-g,\, f-g\rangle\\ &= \langle f,f\rangle - \langle f,g\rangle - \langle g,f\rangle + \langle g,g\rangle\\ &= \|f\|^2 - \langle f,g\rangle - \langle g,f\rangle + \|g\|^2\\ &= 1 - 1 - 1 + 1\\ &= 0 . \end{aligned} \end{equation*}

The definiteness requirement in the definition of an inner product (8.1) states that \(\langle h,h\rangle = 0\) only when \(h = 0\); equivalently, \(\|h\| = 0\) if and only if \(h = 0\) (as noted after 8.5). Applying this with \(h = f-g\), from \(\|f-g\| = 0\) we conclude \(f - g = 0\), that is, \(f = g\).

Note that this argument works over \(\mathbb{R}\) and over \(\mathbb{C}\) alike; the hypothesis \(\langle f,g\rangle = 1\) (rather than merely \(|\langle f,g\rangle| = 1\)) is what forces \(\langle g,f\rangle = 1\) as well. Observe also that the hypotheses say exactly that equality holds in the Cauchy–Schwarz inequality (8.11), since \(|\langle f,g\rangle| = 1 = \|f\|\,\|g\|\); by the equality condition in 8.11 one of \(f,g\) is a scalar multiple of the other. Both are nonzero, since each has norm \(1\); so in either case we may write \(f = \alpha g\) for some scalar \(\alpha\) (if instead \(g = \beta f\), then \(\beta \neq 0\) and \(f = \beta^{-1}g\)). Then \(1 = \langle f,g\rangle = \alpha\langle g,g\rangle = \alpha\|g\|^2 = \alpha\), so \(f = g\) again. Either route establishes the claim.

Problem (8A.13)

Suppose \(f\) and \(g\) are elements of a real inner product space. Prove that

\begin{equation*} \langle f,g\rangle = \frac{\|f+g\|^2 - \|f-g\|^2}{4}. \end{equation*}

Solution

Because the scalar field is \(\mathbb{R}\), conjugate symmetry reads \(\langle g,f\rangle = \langle f,g\rangle\).

Expanding \(\|f+g\|^2\) by additivity in each slot (using 8.3(b) for the second slot),

\begin{equation*} \begin{aligned} \|f+g\|^2 &= \langle f+g,\, f+g\rangle\\ &= \langle f,f\rangle + \langle f,g\rangle + \langle g,f\rangle + \langle g,g\rangle\\ &= \|f\|^2 + \|g\|^2 + 2\langle f,g\rangle . \end{aligned} \end{equation*}

Replacing \(g\) by \(-g\) (and using 8.3(c) with \(\alpha = -1\), so that \(\langle f,-g\rangle = -\langle f,g\rangle\), together with \(\|-g\| = \|g\|\) from 8.6) gives

\begin{equation*} \|f-g\|^2 = \|f\|^2 + \|g\|^2 - 2\langle f,g\rangle . \end{equation*}

Subtracting the second display from the first, the terms \(\|f\|^2 + \|g\|^2\) cancel and we obtain

\begin{equation*} \|f+g\|^2 - \|f-g\|^2 = 4\langle f,g\rangle . \end{equation*}

Dividing both sides by \(4\) yields

\begin{equation*} \langle f,g\rangle = \frac{\|f+g\|^2 - \|f-g\|^2}{4}, \end{equation*}

as desired.

This identity, known as the polarization identity for real inner product spaces, shows that on a real inner product space the inner product is completely determined by the norm it induces: two real inner products on the same vector space that induce the same norm must be equal. (Compare 8.20, which likewise expresses a relation among \(\|f+g\|\), \(\|f-g\|\), \(\|f\|\), and \(\|g\|\).)

Problem (8A.14)

Suppose \(f\) and \(g\) are elements of a complex inner product space. Prove that

\begin{equation*} \langle f,g\rangle = \frac{\|f+g\|^2 - \|f-g\|^2 + \|f+ig\|^2\,i - \|f-ig\|^2\,i}{4}. \end{equation*}

Solution

Here \(\mathbf{F} = \mathbb{C}\), so conjugate symmetry gives \(\langle g,f\rangle = \overline{\langle f,g\rangle}\), and by 8.3(c) we have \(\langle f,\alpha g\rangle = \overline{\alpha}\langle f,g\rangle\) for \(\alpha \in \mathbb{C}\).

Write \(\langle f,g\rangle = x + yi\) with \(x = \operatorname{Re}\langle f,g\rangle\) and \(y = \operatorname{Im}\langle f,g\rangle\), so that \(\langle g,f\rangle = x - yi\).

Step 1: the real part. Expanding as in the proof of 8.15,

\begin{equation*} \begin{aligned} \|f+g\|^2 &= \langle f,f\rangle + \langle f,g\rangle + \langle g,f\rangle + \langle g,g\rangle = \|f\|^2 + \|g\|^2 + 2\operatorname{Re}\langle f,g\rangle,\\ \|f-g\|^2 &= \langle f,f\rangle - \langle f,g\rangle - \langle g,f\rangle + \langle g,g\rangle = \|f\|^2 + \|g\|^2 - 2\operatorname{Re}\langle f,g\rangle, \end{aligned} \end{equation*}

where we used \(\langle f,g\rangle + \langle g,f\rangle = \langle f,g\rangle + \overline{\langle f,g\rangle} = 2\operatorname{Re}\langle f,g\rangle\). Subtracting,

\begin{equation*} \|f+g\|^2 - \|f-g\|^2 = 4\operatorname{Re}\langle f,g\rangle = 4x . \end{equation*}

Step 2: the imaginary part. Apply the first identity of Step 1 with \(g\) replaced by \(ig\). By 8.6, \(\|ig\| = |i|\,\|g\| = \|g\|\), and by 8.3(c),

\begin{equation*} \langle f, ig\rangle = \overline{i}\,\langle f,g\rangle = -i\langle f,g\rangle, \end{equation*}

so that \(2\operatorname{Re}\langle f,ig\rangle = 2\operatorname{Re}\big(-i(x+yi)\big) = 2\operatorname{Re}(y - xi) = 2y\). Hence

\begin{equation*} \|f+ig\|^2 = \|f\|^2 + \|g\|^2 + 2y . \end{equation*}

Similarly, replacing \(g\) by \(-ig\) (again \(\|-ig\| = \|g\|\), and \(\langle f,-ig\rangle = i\langle f,g\rangle\), whose real part is \(-y\)),

\begin{equation*} \|f-ig\|^2 = \|f\|^2 + \|g\|^2 - 2y . \end{equation*}

Subtracting,

\begin{equation*} \|f+ig\|^2 - \|f-ig\|^2 = 4y = 4\operatorname{Im}\langle f,g\rangle . \end{equation*}

Step 3: combine. Multiplying the conclusion of Step 2 by \(i\) and adding the conclusion of Step 1,

\begin{equation*} \|f+g\|^2 - \|f-g\|^2 + \|f+ig\|^2\,i - \|f-ig\|^2\,i = 4x + 4yi = 4\langle f,g\rangle . \end{equation*}

Dividing by \(4\) gives

\begin{equation*} \langle f,g\rangle = \frac{\|f+g\|^2 - \|f-g\|^2 + \|f+ig\|^2\,i - \|f-ig\|^2\,i}{4}, \end{equation*}

as desired.

As in 8A.13, this polarization identity shows that a complex inner product is determined by its associated norm. Note that the complex case genuinely requires the two extra terms: the real formula of 8A.13 recovers only \(\operatorname{Re}\langle f,g\rangle\) when \(\mathbf{F} = \mathbb{C}\), which is not enough to determine \(\langle f,g\rangle\).

Problem (8A.15)

Suppose \(f, g, h\) are elements of an inner product space. Prove that

\begin{equation*} \left\|h - \tfrac{1}{2}(f+g)\right\|^2 = \frac{\|h-f\|^2 + \|h-g\|^2}{2} - \frac{\|f-g\|^2}{4}. \end{equation*}

Solution

Let \(V\) be the inner product space containing \(f\), \(g\), \(h\), and apply the parallelogram equality (8.20) to the two elements

\begin{equation*} u = h - f, \qquad v = h - g \end{equation*}

of \(V\). The parallelogram equality says \(\|u+v\|^2 + \|u-v\|^2 = 2\|u\|^2 + 2\|v\|^2\), which in this case reads

\begin{equation*} \bigl\|2h - (f+g)\bigr\|^2 + \|g - f\|^2 = 2\|h-f\|^2 + 2\|h-g\|^2 , \end{equation*}

because \(u + v = 2h - (f+g)\) and \(u - v = (h-f) - (h-g) = g - f\).

Now \(2h - (f+g) = 2\bigl(h - \tfrac{1}{2}(f+g)\bigr)\), so homogeneity of the norm (8.6) gives

\begin{equation*} \bigl\|2h - (f+g)\bigr\|^2 = 4\left\|h - \tfrac{1}{2}(f+g)\right\|^2 . \end{equation*}

Also \(\|g-f\| = \|-(f-g)\| = \|f-g\|\), again by 8.6. Substituting these two identities into the displayed equation yields

\begin{equation*} 4\left\|h - \tfrac{1}{2}(f+g)\right\|^2 + \|f-g\|^2 = 2\|h-f\|^2 + 2\|h-g\|^2 . \end{equation*}

Solving for the first term and dividing by \(4\) gives

\begin{equation*} \left\|h - \tfrac{1}{2}(f+g)\right\|^2 = \frac{2\|h-f\|^2 + 2\|h-g\|^2 - \|f-g\|^2}{4} = \frac{\|h-f\|^2 + \|h-g\|^2}{2} - \frac{\|f-g\|^2}{4}, \end{equation*}

as desired.

[Geometrically this is Apollonius’s identity: it computes the distance from a point \(h\) to the midpoint of the segment joining \(f\) and \(g\) in terms of the distances from \(h\) to \(f\) and to \(g\) and the length of that segment. Compare Exercise 22 in this section.]

Problem (8A.16)

Prove that a norm satisfying the parallelogram equality comes from an inner product. In other words, show that if \(V\) is a normed vector space whose norm \(\|\cdot\|\) satisfies the parallelogram equality, then there is an inner product \(\langle \cdot, \cdot \rangle\) on \(V\) such that \(\|f\| = \langle f, f\rangle^{1/2}\) for all \(f \in V\).

Solution

Throughout, \(\|\cdot\|\) is a norm on the vector space \(V\) over \(\mathbf{F}\) satisfying

\begin{equation*} \|u+v\|^2 + \|u-v\|^2 = 2\|u\|^2 + 2\|v\|^2 \qquad \text{for all } u, v \in V. \end{equation*}

Exercises 13 and 14 of this section show that if the norm did come from an inner product, then that inner product would necessarily be given by the polarization formulas below; so those formulas are the only candidates, and the work is to verify that they satisfy the axioms in 8.1.

Step 1: the real case.

Assume first that \(\mathbf{F} = \mathbf{R}\), and define

\begin{equation*} \langle f, g\rangle = \frac{\|f+g\|^2 - \|f-g\|^2}{4} \qquad \text{for } f, g \in V. \end{equation*}

(i) Norm and definiteness. Taking \(g = f\) gives \(\langle f, f\rangle = \frac{1}{4}\|2f\|^2 = \|f\|^2\). Hence \(\langle f, f\rangle \ge 0\), and \(\langle f, f\rangle = 0\) if and only if \(\|f\| = 0\), which happens if and only if \(f = 0\). Also \(\|f\| = \langle f, f\rangle^{1/2}\), which is the required relation between the norm and the inner product.

(ii) Symmetry. Because \(\|g - f\| = \|-(f-g)\| = \|f - g\|\), we have \(\langle g, f\rangle = \langle f, g\rangle\) (and every value is real, so conjugate symmetry holds trivially when \(\mathbf{F} = \mathbf{R}\)).

(iii) Additivity in the first slot. Let \(f_1, f_2, g \in V\). Applying the parallelogram equality to the pair \(f_1 + g\), \(f_2\) and then to the pair \(f_1 - g\), \(f_2\) gives

\begin{equation*} \begin{aligned} \|f_1 + f_2 + g\|^2 + \|f_1 - f_2 + g\|^2 &= 2\|f_1+g\|^2 + 2\|f_2\|^2, \\ \|f_1 + f_2 - g\|^2 + \|f_1 - f_2 - g\|^2 &= 2\|f_1-g\|^2 + 2\|f_2\|^2 . \end{aligned} \end{equation*}

Subtracting the second equation from the first and dividing by \(4\) gives

\begin{equation*} \langle f_1 + f_2, g\rangle + \langle f_1 - f_2, g\rangle = \frac{\|f_1+g\|^2 - \|f_1-g\|^2}{2} = 2\langle f_1, g\rangle. \qquad (*) \end{equation*}

Note that \(\langle 0, g\rangle = \frac{1}{4}(\|g\|^2 - \|-g\|^2) = 0\). Thus taking \(f_2 = f_1\) in \((*)\) gives

\begin{equation*} \langle 2f, g\rangle = 2\langle f, g\rangle \qquad \text{for all } f, g \in V. \end{equation*}

Now let \(u, v \in V\) and apply \((*)\) with \(f_1 = \frac{u+v}{2}\) and \(f_2 = \frac{u-v}{2}\), so that \(f_1 + f_2 = u\) and \(f_1 - f_2 = v\):

\begin{equation*} \langle u, g\rangle + \langle v, g\rangle = 2\left\langle \tfrac{u+v}{2}, g\right\rangle = \langle u+v, g\rangle, \end{equation*}

where the last equality uses the doubling identity applied to \(\frac{u+v}{2}\). Thus \(f \mapsto \langle f, g\rangle\) is additive.

(iv) Homogeneity in the first slot. From the definition,

\begin{equation*} \langle -f, g\rangle = \frac{\|{-f}+g\|^2 - \|{-f}-g\|^2}{4} = \frac{\|f-g\|^2 - \|f+g\|^2}{4} = -\langle f, g\rangle . \end{equation*}

Combining this with additivity and induction gives \(\langle nf, g\rangle = n\langle f, g\rangle\) for every \(n \in \mathbf{Z}\) and all \(f, g \in V\). If \(m\) is a positive integer, then applying this with \(n = m\) and \(f\) replaced by \(\frac{1}{m}f\) gives \(\langle f, g\rangle = m\bigl\langle \frac{1}{m}f, g\bigr\rangle\), so \(\bigl\langle \frac{1}{m}f, g\bigr\rangle = \frac{1}{m}\langle f, g\rangle\). Hence

\begin{equation*} \langle tf, g\rangle = t\langle f, g\rangle \qquad \text{for all } t \in \mathbf{Q}, \ f, g \in V. \end{equation*}

Fix \(f, g \in V\). The map \(t \mapsto \|tf + g\|\) is continuous on \(\mathbf{R}\), because

\begin{equation*} \bigl|\,\|tf+g\| - \|sf+g\|\,\bigr| \le \|(t-s)f\| = |t-s|\,\|f\| \end{equation*}

by the triangle inequality; similarly \(t \mapsto \|tf-g\|\) is continuous. Hence \(t \mapsto \langle tf, g\rangle\) is a continuous function of \(t \in \mathbf{R}\), as is \(t \mapsto t\langle f, g\rangle\). These two continuous functions agree on the dense subset \(\mathbf{Q}\) of \(\mathbf{R}\), so they agree on all of \(\mathbf{R}\). Thus \(\langle tf, g\rangle = t\langle f, g\rangle\) for all \(t \in \mathbf{R}\).

By (iii) and (iv), \(f \mapsto \langle f, g\rangle\) is a linear map from \(V\) to \(\mathbf{R}\) for each \(g \in V\). Together with (i) and (ii), all the requirements of 8.1 hold, so \(\langle \cdot, \cdot\rangle\) is an inner product on \(V\) whose associated norm is \(\|\cdot\|\).

Step 2: the complex case.

Now assume \(\mathbf{F} = \mathbf{C}\). Regard \(V\) also as a vector space over \(\mathbf{R}\) (restrict scalar multiplication to real scalars); the norm is unchanged and still satisfies the parallelogram equality. Hence Step 1 applies to this real vector space: the formula

\begin{equation*} \langle f, g\rangle_{\mathbf{R}} = \frac{\|f+g\|^2 - \|f-g\|^2}{4} \end{equation*}

defines a real inner product on \(V\) viewed as a real vector space. In particular \(\langle \cdot, \cdot\rangle_{\mathbf{R}}\) is real valued, symmetric, additive and \(\mathbf{R}\)-homogeneous in each slot, and \(\langle f, f\rangle_{\mathbf{R}} = \|f\|^2\).

Define

\begin{equation*} \langle f, g\rangle = \langle f, g\rangle_{\mathbf{R}} + i\,\langle f, ig\rangle_{\mathbf{R}} = \frac{\|f+g\|^2 - \|f-g\|^2 + i\|f+ig\|^2 - i\|f-ig\|^2}{4} . \end{equation*}

We first record an invariance. For all \(f, g \in V\),

\begin{equation*} \langle if, ig\rangle_{\mathbf{R}} = \frac{\|i(f+g)\|^2 - \|i(f-g)\|^2}{4} = \frac{\|f+g\|^2 - \|f-g\|^2}{4} = \langle f, g\rangle_{\mathbf{R}}, \qquad (\dagger) \end{equation*}

where we used \(\|i u\| = |i|\,\|u\| = \|u\|\) (homogeneity of the norm on the complex normed space \(V\)). Replacing \(g\) by \(ig\) in \((\dagger)\) and using \(i(ig) = -g\) gives

\begin{equation*} \langle f, ig\rangle_{\mathbf{R}} = \langle if, -g\rangle_{\mathbf{R}} = -\langle if, g\rangle_{\mathbf{R}} . \qquad (\ddagger) \end{equation*}

(i) \(\langle f, f\rangle = \|f\|^2\). Indeed \(\langle f, if\rangle_{\mathbf{R}} = \frac{1}{4}\bigl(\|f+if\|^2 - \|f-if\|^2\bigr) = \frac{1}{4}\bigl(|1+i|^2\|f\|^2 - |1-i|^2\|f\|^2\bigr) = \frac{1}{4}(2\|f\|^2 - 2\|f\|^2) = 0\). Hence

\begin{equation*} \langle f, f\rangle = \langle f, f\rangle_{\mathbf{R}} + i\langle f, if\rangle_{\mathbf{R}} = \|f\|^2 . \end{equation*}

Thus \(\langle f, f\rangle \ge 0\), with equality if and only if \(f = 0\), and \(\|f\| = \langle f, f\rangle^{1/2}\).

(ii) Additivity and real homogeneity in the first slot. Both \(f \mapsto \langle f, g\rangle_{\mathbf{R}}\) and \(f \mapsto \langle f, ig\rangle_{\mathbf{R}}\) are additive and \(\mathbf{R}\)-homogeneous by Step 1, hence so is \(f \mapsto \langle f, g\rangle\).

(iii) \(\langle if, g\rangle = i\langle f, g\rangle\). Using the definition, then \((\dagger)\), then \((\ddagger)\):

\begin{equation*} \langle if, g\rangle = \langle if, g\rangle_{\mathbf{R}} + i\langle if, ig\rangle_{\mathbf{R}} = -\langle f, ig\rangle_{\mathbf{R}} + i\langle f, g\rangle_{\mathbf{R}} = i\bigl(\langle f, g\rangle_{\mathbf{R}} + i\langle f, ig\rangle_{\mathbf{R}}\bigr) = i\langle f, g\rangle . \end{equation*}

Combining (ii) and (iii): if \(\alpha = s + it\) with \(s, t \in \mathbf{R}\), then

\begin{equation*} \langle \alpha f, g\rangle = \langle sf, g\rangle + \langle t(if), g\rangle = s\langle f, g\rangle + t\,i\langle f, g\rangle = \alpha \langle f, g\rangle . \end{equation*}

Thus \(f \mapsto \langle f, g\rangle\) is a linear map from \(V\) to \(\mathbf{C}\) for each \(g \in V\).

(iv) Conjugate symmetry. Using the symmetry of \(\langle \cdot, \cdot\rangle_{\mathbf{R}}\) and then \((\ddagger)\),

\begin{equation*} \langle g, f\rangle = \langle g, f\rangle_{\mathbf{R}} + i\langle g, if\rangle_{\mathbf{R}} = \langle f, g\rangle_{\mathbf{R}} + i\langle if, g\rangle_{\mathbf{R}} = \langle f, g\rangle_{\mathbf{R}} - i\langle f, ig\rangle_{\mathbf{R}} = \overline{\langle f, g\rangle}, \end{equation*}

where the last equality holds because \(\langle f, g\rangle_{\mathbf{R}}\) and \(\langle f, ig\rangle_{\mathbf{R}}\) are real numbers.

By (i)-(iv), \(\langle \cdot, \cdot \rangle\) satisfies all the requirements of 8.1, so it is an inner product on \(V\), and by (i) its associated norm is the given norm \(\|\cdot\|\).

In both cases we have produced an inner product on \(V\) with \(\|f\| = \langle f, f\rangle^{1/2}\) for all \(f \in V\), completing the proof.

[Together with 8.20, this shows that a norm comes from an inner product exactly when it satisfies the parallelogram equality. This result is due to Jordan and von Neumann (1935).]

Problem (8A.17)

Let \(\lambda\) denote Lebesgue measure on \([1, \infty)\).

(a) Prove that if \(f \colon [1, \infty) \to [0, \infty)\) is Borel measurable, then

\begin{equation*} \left(\int_1^{\infty} f(x) \, d\lambda(x)\right)^{2} \le \int_1^{\infty} x^2 \bigl(f(x)\bigr)^2 \, d\lambda(x). \end{equation*}

(b) Describe the set of Borel measurable functions \(f \colon [1, \infty) \to [0, \infty)\) such that the inequality in (a) is an equality.

Solution

Define Borel measurable functions \(g, h \colon [1,\infty) \to [0,\infty)\) by

\begin{equation*} g(x) = x f(x), \qquad h(x) = \frac{1}{x}. \end{equation*}

Then \(g(x) h(x) = f(x)\) for every \(x \in [1,\infty)\), and \(\bigl(g(x)\bigr)^2 = x^2 \bigl(f(x)\bigr)^2\).

First we compute \(\|h\|_2\). The functions \(x \mapsto x^{-2}\chi_{[1,n]}(x)\) increase pointwise to \(x \mapsto x^{-2}\) on \([1,\infty)\) as \(n \to \infty\), so the Monotone Convergence Theorem (3.11) gives

\begin{equation*} \int_1^{\infty} \frac{1}{x^2} \, d\lambda(x) = \lim_{n \to \infty} \int_1^{n} \frac{1}{x^2} \, d\lambda(x) = \lim_{n \to \infty}\left(1 - \frac{1}{n}\right) = 1, \end{equation*}

where the middle equality holds because \(x \mapsto x^{-2}\) is continuous on \([1,n]\) and hence Riemann integrable there, so its Lebesgue integral over \([1,n]\) equals its Riemann integral (3.34). Thus \(h \in L^2(\lambda)\) and \(\|h\|_2 = 1\).

(a) If \(\int_1^{\infty} x^2 (f(x))^2 \, d\lambda(x) = \infty\), then the right side of the desired inequality is \(\infty\) and there is nothing to prove. So assume

\begin{equation*} \int_1^{\infty} \bigl(g(x)\bigr)^2 \, d\lambda(x) = \int_1^{\infty} x^2 \bigl(f(x)\bigr)^2 \, d\lambda(x) < \infty, \end{equation*}

which means \(g \in L^2(\lambda)\). Because \(g\) and \(h\) are nonnegative real valued, \(\langle g, h\rangle = \int_1^{\infty} g\overline{h}\,d\lambda = \int_1^{\infty} gh\,d\lambda = \int_1^{\infty} f\,d\lambda\), a nonnegative number, so \(\int_1^{\infty} f\,d\lambda = |\langle g, h\rangle|\). Applying the Cauchy–Schwarz inequality for \(L^2(\lambda)\) (8.14, which is 8.11 applied to the standard inner product on \(L^2(\lambda)\)) to \(g\) and \(h\) now gives

\begin{equation*} \int_1^{\infty} f \, d\lambda = \bigl|\langle g, h\rangle\bigr| \le \|g\|_2 \, \|h\|_2 = \left(\int_1^{\infty} x^2 \bigl(f(x)\bigr)^2 \, d\lambda(x)\right)^{1/2}, \end{equation*}

using \(\|h\|_2 = 1\). In particular \(\int_1^\infty f\,d\lambda < \infty\). Squaring both sides of the inequality above (both sides are nonnegative) gives

\begin{equation*} \left(\int_1^{\infty} f \, d\lambda\right)^{2} \le \int_1^{\infty} x^2 \bigl(f(x)\bigr)^2 \, d\lambda(x), \end{equation*}

as desired.

(b) The claim is that equality holds if and only if

\begin{equation*} \text{either} \quad \int_1^{\infty} f \, d\lambda = \infty, \quad \text{or} \quad f(x) = \frac{c}{x^2} \ \text{ for almost every } x \in [1,\infty), \text{ for some } c \in [0,\infty). \end{equation*}

First we check that each of these conditions gives equality. If \(\int_1^{\infty} f \, d\lambda = \infty\), then the left side of the inequality in (a) equals \(\infty\), and then (a) forces the right side to equal \(\infty\) as well; thus both sides equal \(\infty\) and we have equality. If \(f(x) = c/x^2\) for almost every \(x\), with \(c \in [0,\infty)\), then

\begin{equation*} \int_1^{\infty} f \, d\lambda = c\int_1^{\infty} \frac{1}{x^2}\,d\lambda(x) = c \quad \text{and} \quad \int_1^{\infty} x^2 \bigl(f(x)\bigr)^2 d\lambda(x) = c^2\int_1^{\infty}\frac{1}{x^2}\,d\lambda(x) = c^2 , \end{equation*}

so both sides of the inequality equal \(c^2\).

Conversely, suppose equality holds and \(\int_1^{\infty} f \, d\lambda < \infty\). Then the right side of the inequality in (a) equals \(\bigl(\int_1^{\infty} f\, d\lambda\bigr)^2 < \infty\), so \(g \in L^2(\lambda)\) and the proof of (a) applies. Because \(g\) and \(h\) are nonnegative,

\begin{equation*} \langle g, h\rangle = \int_1^{\infty} g h \, d\lambda = \int_1^{\infty} f \, d\lambda , \end{equation*}

and our equality says precisely that \(|\langle g, h\rangle| = \|g\|_2 \|h\|_2\). By the condition for equality in the Cauchy–Schwarz inequality (8.11), one of \(g\), \(h\) is a scalar multiple of the other as elements of \(L^2(\lambda)\). Since \(h \ne 0\) in \(L^2(\lambda)\), in either case there is a scalar \(\alpha\) with \(g = \alpha h\) in \(L^2(\lambda)\): this is immediate if \(g\) is a multiple of \(h\), and if instead \(h = \beta g\) then \(\beta \ne 0\) (as \(h \ne 0\)) and \(g = \beta^{-1} h\). Moreover

\begin{equation*} \alpha = \alpha \|h\|_2^2 = \langle \alpha h, h\rangle = \langle g, h\rangle = \int_1^{\infty} f\, d\lambda \ge 0 . \end{equation*}

Now \(g = \alpha h\) in \(L^2(\lambda)\) means \(x f(x) = \alpha/x\) for almost every \(x \in [1,\infty)\), that is,

\begin{equation*} f(x) = \frac{\alpha}{x^2} \quad \text{for almost every } x \in [1,\infty), \end{equation*}

with \(\alpha \in [0,\infty)\). This proves the claim.

Thus, ignoring the degenerate family of functions with divergent integral (for which both sides are \(\infty\)), equality holds exactly for the functions that agree almost everywhere with a nonnegative constant multiple of \(1/x^2\).

Problem (8A.18)

Suppose \(\mu\) is a measure. For \(f, g \in L^2(\mu)\), define \(\langle f, g\rangle\) by

\begin{equation*} \langle f, g\rangle = \int f \overline{g} \, d\mu . \end{equation*}

(a) Using the inequality

\begin{equation*} \bigl|f(x)\,\overline{g(x)}\bigr| \le \tfrac{1}{2}\bigl(|f(x)|^2 + |g(x)|^2\bigr), \end{equation*}

verify that the integral above makes sense and the map sending \(f, g\) to \(\langle f, g\rangle\) defines an inner product on \(L^2(\mu)\) (without using Hölder’s inequality).

(b) Show that the Cauchy–Schwarz inequality implies that

\begin{equation*} \|fg\|_1 \le \|f\|_2 \, \|g\|_2 \end{equation*}

for all \(f, g \in L^2(\mu)\) (again, without using Hölder’s inequality).

Solution

Let \((X, \mathcal{S}, \mu)\) be the measure space in question.

(a) First note the stated inequality. If \(a, b \in [0,\infty)\), then \(0 \le (a-b)^2 = a^2 - 2ab + b^2\), so \(ab \le \frac{1}{2}(a^2 + b^2)\). Taking \(a = |f(x)|\) and \(b = |g(x)| = \bigl|\overline{g(x)}\bigr|\) gives

\begin{equation*} \bigl|f(x)\overline{g(x)}\bigr| = |f(x)| \, |g(x)| \le \tfrac{1}{2}\bigl(|f(x)|^2 + |g(x)|^2\bigr) \qquad \text{for all } x \in X . \end{equation*}

The integral makes sense. Suppose \(f, g \in L^2(\mu)\). The function \(f\overline{g}\) is \(\mathcal{S}\)-measurable (the conjugate of an \(\mathcal{S}\)-measurable complex function is \(\mathcal{S}\)-measurable, and products of \(\mathcal{S}\)-measurable functions are \(\mathcal{S}\)-measurable), and \(|f\overline{g}|\) is \(\mathcal{S}\)-measurable by 6.20. Because \(f, g \in L^2(\mu)\), the nonnegative functions \(|f|^2\) and \(|g|^2\) have finite integrals. Hence, by the displayed inequality and because integration is order preserving (3.8) and additive,

\begin{equation*} \int |f\overline{g}| \, d\mu \le \tfrac{1}{2}\int \bigl(|f|^2 + |g|^2\bigr) \, d\mu = \tfrac{1}{2}\bigl(\|f\|_2^2 + \|g\|_2^2\bigr) < \infty . \end{equation*}

Thus \(\int f \overline{g}\, d\mu\) is a well-defined element of \(\mathbf{F}\) (by 6.21 in the case \(\mathbf{F} = \mathbf{C}\); the real case is 3.18). Note that no use has been made of Hölder’s inequality.

The definition is also independent of the representatives chosen: if \(f = f^{\prime}\) almost everywhere and \(g = g^{\prime}\) almost everywhere, then \(f\overline{g} = f^{\prime}\overline{g^{\prime}}\) almost everywhere, so the two integrals are equal. Hence \(\langle \cdot, \cdot\rangle\) is a well-defined map from \(L^2(\mu) \times L^2(\mu)\) to \(\mathbf{F}\).

It is an inner product. We check the requirements in 8.1.

Linearity in the first slot: if \(f_1, f_2, g \in L^2(\mu)\) and \(\alpha \in \mathbf{F}\), then \(f_1\overline{g}\), \(f_2\overline{g}\), and \((\alpha f_1 + f_2)\overline{g} = \alpha f_1 \overline{g} + f_2 \overline{g}\) all have finite integrals of their absolute values by the previous paragraph, so homogeneity and additivity of integration (3.20 and 3.21, applied to real and imaginary parts as in the discussion following 6.21) give

\begin{equation*} \langle \alpha f_1 + f_2, g\rangle = \int \bigl(\alpha f_1\overline{g} + f_2\overline{g}\bigr) \, d\mu = \alpha \int f_1 \overline{g}\,d\mu + \int f_2\overline{g}\, d\mu = \alpha\langle f_1, g\rangle + \langle f_2, g\rangle . \end{equation*}

Thus \(f \mapsto \langle f, g\rangle\) is a linear map from \(L^2(\mu)\) to \(\mathbf{F}\) for each \(g \in L^2(\mu)\).

Conjugate symmetry: because \(\overline{f\overline{g}} = \overline{f}\,g = g\overline{f}\) and the integral of the complex conjugate is the conjugate of the integral (6.24),

\begin{equation*} \langle g, f\rangle = \int g\overline{f} \, d\mu = \int \overline{f \overline{g}} \, d\mu = \overline{\int f\overline{g}\,d\mu} = \overline{\langle f, g\rangle} . \end{equation*}

Positivity: \(\langle f, f\rangle = \int f\overline{f}\,d\mu = \int |f|^2 \, d\mu = \|f\|_2^2 \ge 0\).

Definiteness: suppose \(\langle f, f\rangle = 0\), that is, \(\int |f|^2 d\mu = 0\). For each \(n \in \mathbf{Z}^+\) let \(E_n = \{x \in X : |f(x)|^2 > 1/n\}\), which is in \(\mathcal{S}\). Then \(\frac{1}{n}\chi_{E_n} \le |f|^2\) pointwise, so \(\frac{1}{n}\mu(E_n) \le \int |f|^2 d\mu = 0\) by 3.8, giving \(\mu(E_n) = 0\). Because \(\{x \in X : f(x) \ne 0\} = \bigcup_{n=1}^{\infty} E_n\), countable subadditivity of measures (2.58) gives \(\mu\bigl(\{x : f(x) \ne 0\}\bigr) = 0\). Thus \(f = 0\) almost everywhere, which says exactly that \(f\) is the \(0\) element of \(L^2(\mu)\). (This is the reason the inner product is defined on \(L^2(\mu)\), whose elements are equivalence classes, rather than on \(\mathcal{L}^2(\mu)\); see 8.2.)

Hence \(\langle \cdot, \cdot\rangle\) is an inner product on \(L^2(\mu)\), and by the positivity computation its associated norm (8.4) is the usual norm \(\|\cdot\|_2\) on \(L^2(\mu)\).

(b) Suppose \(f, g \in L^2(\mu)\). Then \(|f|, |g| \in L^2(\mu)\), because \(\bigl||f|\bigr|^2 = |f|^2\) and \(\bigl||g|\bigr|^2 = |g|^2\) have finite integrals. Applying the Cauchy–Schwarz inequality (8.11) in the inner product space \(L^2(\mu)\) constructed in part (a) to the elements \(|f|\) and \(|g|\), and using that \(|g|\) is real valued (so that \(\overline{|g|} = |g|\)), we get

\begin{equation*} \int |f| \, |g| \, d\mu = \bigl|\bigl\langle |f|, |g| \bigr\rangle\bigr| \le \bigl\| |f| \bigr\| \, \bigl\| |g| \bigr\| = \left(\int |f|^2 d\mu\right)^{1/2}\left(\int |g|^2 d\mu\right)^{1/2} = \|f\|_2 \, \|g\|_2 . \end{equation*}

Because \(|f||g| = |fg|\), the left side equals \(\int |fg| \, d\mu = \|fg\|_1\) (which is finite by part (a), applied with \(g\) replaced by \(\overline{g}\), so \(fg \in L^1(\mu)\)). Therefore

\begin{equation*} \|fg\|_1 \le \|f\|_2\, \|g\|_2 , \end{equation*}

as desired. Note that the whole chain — the definition of the inner product, the Cauchy–Schwarz inequality (whose proof used only 8.9 and 8.10), and this last step — never invoked Hölder’s inequality. Thus this exercise provides a proof of Hölder’s inequality (7.9) in the case \(p = p^{\prime} = 2\) that is independent of 7.9; compare 8.14.

Problem (8A.19)

Suppose \(V_1, \ldots, V_m\) are inner product spaces. Show that the equation

\begin{equation*} \bigl\langle (f_1, \ldots, f_m), (g_1, \ldots, g_m)\bigr\rangle = \langle f_1, g_1\rangle + \cdots + \langle f_m, g_m\rangle \end{equation*}

defines an inner product on \(V_1 \times \cdots \times V_m\).

[Each of the inner product spaces \(V_1, \ldots, V_m\) may have a different inner product, even though the same inner product notation is used on all these spaces.]

Solution

All the spaces \(V_1, \ldots, V_m\) are vector spaces over the same field \(\mathbf{F}\). Recall that \(V = V_1 \times \cdots \times V_m\) is a vector space over \(\mathbf{F}\) with the coordinatewise operations

\begin{equation*} (f_1, \ldots, f_m) + (g_1, \ldots, g_m) = (f_1 + g_1, \ldots, f_m + g_m), \qquad \alpha(f_1, \ldots, f_m) = (\alpha f_1, \ldots, \alpha f_m) \end{equation*}

for \(\alpha \in \mathbf{F}\); its additive identity is \((0, \ldots, 0)\). For \(f = (f_1,\ldots,f_m)\) and \(g = (g_1,\ldots,g_m)\) in \(V\), define

\begin{equation*} \langle f, g\rangle = \sum_{k=1}^{m} \langle f_k, g_k\rangle , \end{equation*}

where the \(k\)-th summand uses the inner product of \(V_k\). This is a well-defined element of \(\mathbf{F}\). We verify the four requirements of 8.1.

Positivity. For \(f = (f_1, \ldots, f_m) \in V\),

\begin{equation*} \langle f, f\rangle = \sum_{k=1}^{m} \langle f_k, f_k\rangle = \sum_{k=1}^{m} \|f_k\|^2 , \end{equation*}

where \(\|\cdot\|\) denotes the norm of the relevant \(V_k\) (8.4). Each term is a nonnegative real number, so \(\langle f, f\rangle \ge 0\).

Definiteness. Suppose \(\langle f, f\rangle = 0\). Because \(\sum_{k=1}^m \|f_k\|^2 = 0\) and each \(\|f_k\|^2 \ge 0\), we get \(\|f_k\|^2 = 0\), hence \(\langle f_k, f_k \rangle = 0\), hence \(f_k = 0\) for each \(k\) (by the definiteness requirement in \(V_k\)). Thus \(f = (0,\ldots,0)\), the \(0\) element of \(V\). Conversely \(\langle 0, 0\rangle = 0\) by 8.3(a) applied in each \(V_k\).

Linearity in the first slot. Fix \(g = (g_1, \ldots, g_m) \in V\). If \(f = (f_1,\ldots,f_m)\) and \(h = (h_1,\ldots,h_m)\) are in \(V\) and \(\alpha \in \mathbf{F}\), then \(\alpha f + h = (\alpha f_1 + h_1, \ldots, \alpha f_m + h_m)\), so

\begin{equation*} \begin{aligned} \langle \alpha f + h, g\rangle &= \sum_{k=1}^{m} \langle \alpha f_k + h_k, g_k\rangle \\ &= \sum_{k=1}^{m} \bigl(\alpha \langle f_k, g_k\rangle + \langle h_k, g_k\rangle\bigr) \\ &= \alpha \sum_{k=1}^{m}\langle f_k, g_k\rangle + \sum_{k=1}^{m}\langle h_k, g_k\rangle = \alpha\langle f, g\rangle + \langle h, g\rangle , \end{aligned} \end{equation*}

where the second equality holds because \(u \mapsto \langle u, g_k\rangle\) is a linear map on \(V_k\) for each \(k\). Thus \(f \mapsto \langle f, g\rangle\) is a linear map from \(V\) to \(\mathbf{F}\).

Conjugate symmetry. Using conjugate symmetry in each \(V_k\) and the additivity of complex conjugation (6.24),

\begin{equation*} \langle g, f\rangle = \sum_{k=1}^{m}\langle g_k, f_k\rangle = \sum_{k=1}^{m}\overline{\langle f_k, g_k\rangle} = \overline{\sum_{k=1}^{m}\langle f_k, g_k\rangle} = \overline{\langle f, g\rangle} . \end{equation*}

All four requirements hold, so \(\langle \cdot, \cdot\rangle\) is an inner product on \(V_1 \times \cdots \times V_m\). Its associated norm is

\begin{equation*} \bigl\|(f_1, \ldots, f_m)\bigr\| = \bigl(\|f_1\|^2 + \cdots + \|f_m\|^2\bigr)^{1/2}, \end{equation*}

as the positivity computation shows.

Problem (8A.20)

Suppose \(V\) is an inner product space. Make \(V \times V\) an inner product space as in the exercise above. Prove that the function that takes an ordered pair \((f, g) \in V \times V\) to the inner product \(\langle f, g\rangle \in \mathbf{F}\) is a continuous function from \(V \times V\) to \(\mathbf{F}\).

Solution

By Exercise 19, \(V \times V\) is an inner product space with

\begin{equation*} \bigl\langle (f_1, f_2), (g_1, g_2)\bigr\rangle = \langle f_1, g_1\rangle + \langle f_2, g_2\rangle , \end{equation*}

and its associated norm is \(\|(f_1,f_2)\| = \bigl(\|f_1\|^2 + \|f_2\|^2\bigr)^{1/2}\). Because \(\|\cdot\|\) is a norm (8.19), \(V \times V\) is a metric space with the metric \(d\bigl((f_1,f_2),(g_1,g_2)\bigr) = \|(f_1 - g_1, f_2 - g_2)\|\), and \(\mathbf{F}\) is a metric space with the metric coming from the absolute value. Continuity of

\begin{equation*} \Phi \colon V \times V \to \mathbf{F}, \qquad \Phi(f, g) = \langle f, g\rangle \end{equation*}

means continuity with respect to these metrics.

Note first that for \((u, v), (f,g) \in V \times V\) we have

\begin{equation*} \|u - f\| \le \bigl(\|u-f\|^2 + \|v-g\|^2\bigr)^{1/2} = d\bigl((u,v),(f,g)\bigr), \end{equation*}

and likewise \(\|v - g\| \le d\bigl((u,v),(f,g)\bigr)\).

Now fix \((f, g) \in V \times V\) and let \(\varepsilon > 0\). Set

\begin{equation*} \delta = \min\left\{1, \ \frac{\varepsilon}{\|f\| + \|g\| + 1}\right\} > 0 . \end{equation*}

Suppose \((u, v) \in V \times V\) satisfies \(d\bigl((u,v),(f,g)\bigr) < \delta\); by the previous paragraph, \(\|u - f\| < \delta\) and \(\|v - g\| < \delta\). Write

\begin{equation*} \langle u, v\rangle - \langle f, g\rangle = \langle u - f, v\rangle + \langle f, v - g\rangle , \end{equation*}

which holds by linearity in the first slot and by 8.3(b). Hence the triangle inequality in \(\mathbf{F}\) and the Cauchy–Schwarz inequality (8.11) give

\begin{equation*} \bigl|\langle u, v\rangle - \langle f, g\rangle\bigr| \le \|u - f\| \, \|v\| + \|f\| \, \|v - g\| . \end{equation*}

Because \(\|v\| \le \|g\| + \|v - g\| < \|g\| + \delta \le \|g\| + 1\) (using the triangle inequality 8.15 and \(\delta \le 1\)), we conclude

\begin{equation*} \bigl|\langle u, v\rangle - \langle f, g\rangle\bigr| < \delta\bigl(\|g\| + 1\bigr) + \|f\|\,\delta = \delta\bigl(\|f\| + \|g\| + 1\bigr) \le \varepsilon . \end{equation*}

Thus for every \(\varepsilon > 0\) there is \(\delta > 0\) such that \(d\bigl((u,v),(f,g)\bigr) < \delta\) implies \(|\Phi(u,v) - \Phi(f,g)| < \varepsilon\). Hence \(\Phi\) is continuous at \((f,g)\), and because \((f,g) \in V \times V\) was arbitrary, \(\Phi\) is continuous on \(V \times V\).

Problem (8A.21)

Suppose \(1 \le p \le \infty\).

(a) Show the norm on \(\ell^p\) comes from an inner product if and only if \(p = 2\).

(b) Show the norm on \(L^p(\mathbf{R})\) comes from an inner product if and only if \(p = 2\).

Solution

Say that a norm \(\|\cdot\|\) on a vector space \(W\) comes from an inner product if there is an inner product \(\langle\cdot,\cdot\rangle\) on \(W\) with \(\|w\| = \langle w, w\rangle^{1/2}\) for all \(w \in W\). The key tool is the parallelogram equality 8.20: if \(\|\cdot\|\) comes from an inner product, then

\begin{equation*} \|u+v\|^2 + \|u-v\|^2 = 2\|u\|^2 + 2\|v\|^2 \qquad \text{for all } u, v \in W. \qquad (\mathrm{P}) \end{equation*}

So to rule out \(p \ne 2\) it suffices to exhibit two elements for which (P) fails.

(a) If \(p = 2\), then the norm \(\|\cdot\|_2\) on \(\ell^2\) is the norm associated with the standard inner product \(\bigl\langle (a_1,a_2,\ldots),(b_1,b_2,\ldots)\bigr\rangle = \sum_{k=1}^{\infty} a_k \overline{b_k}\) on \(\ell^2\); see 8.2 and 8.5. Thus the norm on \(\ell^2\) comes from an inner product.

Conversely, suppose \(1 \le p \le \infty\) and the norm on \(\ell^p\) comes from an inner product. Let

\begin{equation*} u = (1, 0, 0, \ldots), \qquad v = (0, 1, 0, 0, \ldots), \end{equation*}

both of which are in \(\ell^p\). Then \(\|u\|_p = \|v\|_p = 1\), and

\begin{equation*} u + v = (1, 1, 0, 0, \ldots), \qquad u - v = (1, -1, 0, 0, \ldots). \end{equation*}

If \(p < \infty\), then \(\|u+v\|_p = \|u-v\|_p = (1^p + 1^p)^{1/p} = 2^{1/p}\), so (P) becomes

\begin{equation*} 2 \cdot 2^{2/p} = 2 + 2 = 4, \quad \text{that is,} \quad 2^{1 + 2/p} = 2^2 . \end{equation*}

Because \(t \mapsto 2^t\) is injective on \(\mathbf{R}\), this forces \(1 + \frac{2}{p} = 2\), hence \(p = 2\).

If \(p = \infty\), then \(\|u+v\|_{\infty} = \|u-v\|_{\infty} = 1\), so the left side of (P) equals \(1 + 1 = 2\) while the right side equals \(2 + 2 = 4\). Thus (P) fails and the norm on \(\ell^{\infty}\) does not come from an inner product.

Therefore the norm on \(\ell^p\) comes from an inner product if and only if \(p = 2\).

(b) Here \(L^p(\mathbf{R})\) means \(L^p(\lambda)\), where \(\lambda\) is Lebesgue measure on \(\mathbf{R}\). If \(p = 2\), then the norm \(\|\cdot\|_2\) on \(L^2(\lambda)\) is the norm associated with the standard inner product \(\langle f, g\rangle = \int f\overline{g}\,d\lambda\) (see 8.2 and 8.5, or Exercise 18 of this section), so it comes from an inner product.

Conversely, suppose \(1 \le p \le \infty\) and the norm on \(L^p(\mathbf{R})\) comes from an inner product. Let

\begin{equation*} u = \chi_{(0,1)}, \qquad v = \chi_{(1,2)} . \end{equation*}

Both are bounded Borel measurable functions vanishing outside a set of finite measure, so \(u, v \in L^p(\mathbf{R})\) for every \(p \in [1,\infty]\). The intervals \((0,1)\) and \((1,2)\) are disjoint and each has measure \(1\), and \((0,1) \cup (1,2) = (0,2) \setminus \{1\}\). Because \(\lambda(\{1\}) = 0\), the functions \(u+v\) and \(|u-v|\) both agree with \(\chi_{(0,2)}\) except at the single point \(1\), so

\begin{equation*} u + v = \chi_{(0,2)} \quad \text{and} \quad |u - v| = \bigl|\chi_{(0,1)} - \chi_{(1,2)}\bigr| = \chi_{(0,2)} \qquad \text{almost everywhere.} \end{equation*}

Because \(L^p(\mathbf{R})\) consists of equivalence classes of functions that agree almost everywhere, and both \(\|\cdot\|_p\) for \(p < \infty\) and the essential supremum norm depend only on the almost-everywhere equivalence class, these identities may be used freely in computing norms.

Suppose first that \(p < \infty\). Then

\begin{equation*} \|u\|_p = \|v\|_p = 1^{1/p} = 1, \qquad \|u+v\|_p = \|u-v\|_p = \left(\int \chi_{(0,2)}\,d\lambda\right)^{1/p} = 2^{1/p} , \end{equation*}

so exactly as in part (a) the equality (P) becomes \(2^{1+2/p} = 4\), which forces \(p = 2\).

If \(p = \infty\), then \(\|u\|_{\infty} = \|v\|_{\infty} = 1\) and \(\|u+v\|_{\infty} = \|u - v\|_{\infty} = 1\) (the essential supremum of \(\chi_{(0,2)}\) is \(1\)), so the left side of (P) equals \(2\) and the right side equals \(4\). Thus (P) fails and the norm on \(L^{\infty}(\mathbf{R})\) does not come from an inner product.

Therefore the norm on \(L^p(\mathbf{R})\) comes from an inner product if and only if \(p = 2\).

[The same two-disjoint-bumps computation shows that for a general measure \(\mu\), the norm on \(L^p(\mu)\) fails to come from an inner product for \(p \ne 2\), provided \(\mu\) has two disjoint measurable sets of equal finite positive measure. Combined with Exercise 16, we could equally well have argued that the parallelogram equality is not merely necessary but sufficient; here only necessity was needed.]

Problem (8A.22)

Use inner products to prove Apollonius’s identity:

In a triangle with sides of length \(a\), \(b\), and \(c\), let \(d\) be the length of the line segment from the midpoint of the side of length \(c\) to the opposite vertex. Then

\begin{equation*} a^2 + b^2 = \tfrac{1}{2}c^2 + 2d^2 . \end{equation*}

Solution

Work in \(\mathbb{R}^2\) with the standard inner product \(\langle x, y \rangle = x_1 y_1 + x_2 y_2\), whose associated norm is the usual Euclidean norm (see 8.5); thus the distance between two points \(x, y \in \mathbb{R}^2\) equals \(\lVert x - y \rVert\). (Nothing below uses the dimension: the same computation proves the identity for a triangle in any real or complex inner product space.)

Let the triangle have vertices \(P\), \(Q\), \(R\), labelled so that the side of length \(c\) is the segment from \(P\) to \(Q\); then \(R\) is the vertex opposite that side, and the two sides meeting at \(R\) have lengths \(a\) and \(b\). Say

\begin{equation*} a = \lVert Q - R \rVert, \qquad b = \lVert P - R \rVert, \qquad c = \lVert P - Q \rVert . \end{equation*}

(The identity to be proved is symmetric in \(a\) and \(b\), so this choice of which of the two sides at \(R\) is called \(a\) is harmless.)

Translate the picture so that the computation is about two vectors based at \(R\): set

\begin{equation*} f = P - R, \qquad g = Q - R . \end{equation*}

Then

\begin{equation*} \lVert f \rVert = b, \qquad \lVert g \rVert = a, \qquad \lVert f - g \rVert = \lVert P - Q \rVert = c , \end{equation*}

where the last equality holds because \(f - g = (P - R) - (Q - R) = P - Q\).

The midpoint of the side of length \(c\) is \(M = \frac{1}{2}(P + Q)\), and \(d\) is the distance from \(M\) to the opposite vertex \(R\). Since

\begin{equation*} M - R = \tfrac{1}{2}(P + Q) - R = \tfrac{1}{2}\bigl( (P - R) + (Q - R) \bigr) = \tfrac{1}{2}(f + g), \end{equation*}

homogeneity of the norm (8.6, with \(\alpha = \tfrac{1}{2}\)) gives

\begin{equation*} d = \lVert M - R \rVert = \tfrac{1}{2}\lVert f + g \rVert, \qquad \text{so} \qquad \lVert f + g \rVert = 2d . \end{equation*}

Now apply the parallelogram equality 8.20 to \(f\) and \(g\):

\begin{equation*} \lVert f + g \rVert^2 + \lVert f - g \rVert^2 = 2\lVert f \rVert^2 + 2\lVert g \rVert^2 . \end{equation*}

Substituting the four quantities computed above turns this into

\begin{equation*} (2d)^2 + c^2 = 2b^2 + 2a^2, \qquad \text{that is,} \qquad 4d^2 + c^2 = 2a^2 + 2b^2 . \end{equation*}

Dividing both sides by \(2\) gives

\begin{equation*} 2d^2 + \tfrac{1}{2}c^2 = a^2 + b^2 , \end{equation*}

which is Apollonius’s identity.

For completeness, here is the one-line proof of 8.20 in this concrete setting, which shows exactly where the inner product enters. Expanding by additivity and homogeneity of the inner product in each slot,

\begin{equation*} \begin{aligned} \lVert f + g \rVert^2 &= \langle f + g, f + g \rangle = \lVert f \rVert^2 + \langle f, g \rangle + \langle g, f \rangle + \lVert g \rVert^2, \\ \lVert f - g \rVert^2 &= \langle f - g, f - g \rangle = \lVert f \rVert^2 - \langle f, g \rangle - \langle g, f \rangle + \lVert g \rVert^2 . \end{aligned} \end{equation*}

Adding these two equations cancels the cross terms \(\langle f, g \rangle\) and \(\langle g, f \rangle\) and leaves \(2\lVert f \rVert^2 + 2\lVert g \rVert^2\), as claimed. The cancellation of the cross terms is the whole content of the proof: it is why Apollonius’s identity holds for every triangle, with no hypothesis relating the angle at \(R\) to anything else.

Exercises 8B

Problem (8B.1)

Show that each of the inner product spaces in Example 8.23 is not a Hilbert space.

Solution

Example 8.23 lists two inner product spaces: the space \(\ell^1\) with \(\langle (a_1, a_2, \ldots), (b_1, b_2, \ldots)\rangle = \sum_{k=1}^{\infty} a_k \overline{b_k}\), and the space \(C([0,1])\) of continuous \(\mathbf{F}\)-valued functions on \([0,1]\) with \(\langle f, g\rangle = \int_0^1 f \overline{g}\). In each case the associated norm is \(\|f\| = \langle f, f\rangle^{1/2}\) (a norm by 8.19), and by 8.21 the space fails to be a Hilbert space exactly when that norm is not complete. So in each case we exhibit a Cauchy sequence with no limit in the space.

First note a general remark used twice below: if \(\|h_n - h\| \to 0\) then \(\langle h_n - h, h_n - h\rangle \to 0\), so any quantity dominated by \(\|h_n - h\|\) tends to \(0\).

Part 1: the space \(\ell^1\).

Here the norm determined by the inner product is \(\|(a_1, a_2, \ldots)\| = \left(\sum_{k=1}^{\infty} |a_k|^2\right)^{1/2}\), the \(\ell^2\)-norm, restricted to \(\ell^1\). (The inner product is well defined on \(\ell^1\): if \(a, b \in \ell^1\) then \(a, b\) are bounded sequences, so \(\sum_k |a_k \overline{b_k}| \le \left(\sup_j |b_j|\right)\sum_k |a_k| < \infty\).)

For \(n \in \mathbf{Z}^+\) let

\begin{equation*} a^{(n)} = \Bigl(1, \tfrac12, \tfrac13, \ldots, \tfrac1n, 0, 0, \ldots\Bigr). \end{equation*}

Each \(a^{(n)}\) has only finitely many nonzero coordinates, so \(a^{(n)} \in \ell^1\). If \(m < n\) then

\begin{equation*} \|a^{(n)} - a^{(m)}\|^2 = \sum_{k=m+1}^{n} \frac{1}{k^2} \le \sum_{k=m+1}^{\infty} \frac{1}{k^2}, \end{equation*}

which tends to \(0\) as \(m \to \infty\) because \(\sum_{k=1}^{\infty} 1/k^2 < \infty\). Hence \(a^{(1)}, a^{(2)}, \ldots\) is a Cauchy sequence in \(\ell^1\) with respect to \(\|\cdot\|\).

Suppose \(b = (b_1, b_2, \ldots) \in \ell^1\) and \(\|a^{(n)} - b\| \to 0\). For each fixed \(j\) we have \(|a^{(n)}_j - b_j| \le \|a^{(n)} - b\|\), so \(a^{(n)}_j \to b_j\). But \(a^{(n)}_j = 1/j\) for all \(n \ge j\), so \(b_j = 1/j\) for every \(j \in \mathbf{Z}^+\). Then \(\sum_{j=1}^{\infty} |b_j| = \sum_{j=1}^{\infty} 1/j = \infty\), contradicting \(b \in \ell^1\). Thus the Cauchy sequence \(a^{(1)}, a^{(2)}, \ldots\) has no limit in \(\ell^1\), so this norm is not complete on \(\ell^1\) and \(\ell^1\) with this inner product is not a Hilbert space.

Part 2: the space \(C([0,1])\).

Here the norm determined by the inner product is \(\|f\| = \left(\int_0^1 |f|^2\right)^{1/2}\), where the integral is with respect to Lebesgue measure (equivalently, the Riemann integral, since the integrands are continuous).

For \(n \in \mathbf{Z}^+\) with \(n \ge 3\) define \(f_n \in C([0,1])\) by

\begin{equation*} f_n(x) = \begin{cases} 0 & \text{if } 0 \le x \le \frac12, \\ n\left(x - \frac12\right) & \text{if } \frac12 \le x \le \frac12 + \frac1n, \\ 1 & \text{if } \frac12 + \frac1n \le x \le 1. \end{cases} \end{equation*}

Each \(f_n\) is continuous (the three formulas agree at the two junction points) and \(0 \le f_n \le 1\).

If \(3 \le m \le n\) then \(f_n(x) = f_m(x)\) for \(x \le \frac12\) and for \(x \ge \frac12 + \frac1m\), and \(|f_n - f_m| \le 1\) everywhere; hence

\begin{equation*} \|f_n - f_m\|^2 = \int_{1/2}^{1/2 + 1/m} |f_n - f_m|^2 \le \frac{1}{m}. \end{equation*}

Therefore \(f_3, f_4, \ldots\) is a Cauchy sequence in \(C([0,1])\) with respect to \(\|\cdot\|\).

Suppose \(f \in C([0,1])\) and \(\|f_n - f\| \to 0\). Since \(f_n = 0\) on \([0, \frac12]\) for every \(n\),

\begin{equation*} \int_0^{1/2} |f|^2 = \int_0^{1/2} |f_n - f|^2 \le \|f_n - f\|^2 \to 0, \end{equation*}

so \(\int_0^{1/2} |f|^2 = 0\). Because \(|f|^2\) is continuous and nonnegative, this forces \(f(x) = 0\) for all \(x \in [0, \frac12]\): indeed if \(|f(x_0)| > 0\) for some \(x_0 \in [0,\frac12]\), then by continuity \(|f|^2 > \frac12 |f(x_0)|^2\) on a subinterval of \([0,\frac12]\) of positive length, making the integral positive. In particular \(f(\frac12) = 0\).

Now fix \(a\) with \(\frac12 < a < 1\). For all \(n\) large enough that \(\frac12 + \frac1n < a\) we have \(f_n = 1\) on \([a, 1]\), so

\begin{equation*} \int_a^1 |1 - f|^2 = \int_a^1 |f_n - f|^2 \le \|f_n - f\|^2 \to 0. \end{equation*}

Hence \(\int_a^1 |1 - f|^2 = 0\), and the same continuity argument gives \(f(x) = 1\) for all \(x \in [a,1]\). Since \(a \in (\frac12, 1)\) was arbitrary, \(f = 1\) on \((\frac12, 1]\), and continuity of \(f\) at \(\frac12\) gives \(f(\frac12) = \lim_{x \downarrow 1/2} f(x) = 1\).

This contradicts \(f(\frac12) = 0\). So the Cauchy sequence \(f_3, f_4, \ldots\) has no limit in \(C([0,1])\), the norm is not complete on \(C([0,1])\), and \(C([0,1])\) with this inner product is not a Hilbert space.

Problem (8B.2)

Prove or disprove: The inner product space in Exercise 1 in Section 8A is a Hilbert space.

Solution

The statement is false; we disprove it.

Recall the space from Exercise 1 in Section 8A: \(V\) is the vector space of bounded continuous functions from \(\mathbf{R}\) to \(\mathbf{F}\), a list \(r_1, r_2, \ldots\) of the rational numbers is fixed, and

\begin{equation*} \langle f, g\rangle = \sum_{k=1}^{\infty} \frac{f(r_k)\overline{g(r_k)}}{2^k} \end{equation*}

for \(f, g \in V\). That exercise shows \(\langle \cdot, \cdot\rangle\) is an inner product on \(V\) (positive definiteness uses that \(\mathbf{Q}\) is dense in \(\mathbf{R}\) and \(f\) is continuous: if \(f(r_k) = 0\) for every \(k\) then \(f = 0\)). The associated norm is

\begin{equation*} \|f\| = \Bigl(\sum_{k=1}^{\infty} \frac{|f(r_k)|^2}{2^k}\Bigr)^{1/2}. \end{equation*}

We show this norm is not complete on \(V\); by 8.21, \(V\) is then not a Hilbert space.

Two elementary observations.

First, for each fixed \(j \in \mathbf{Z}^+\) and all \(f, g \in V\),

\begin{equation*} \frac{|f(r_j) - g(r_j)|}{2^{j/2}} \le \|f - g\|, \end{equation*}

since the left side squared is one term of the series defining \(\|f-g\|^2\) and all terms are nonnegative. So norm convergence forces convergence of the values at each \(r_j\).

Second, \(\|f\| \le \sup_{\mathbf{R}} |f|\), because \(\sum_{k=1}^{\infty} 2^{-k} = 1\).

Now define, for \(n \in \mathbf{Z}^+\), the function \(\varphi_n : \mathbf{R} \to \mathbf{F}\) by

\begin{equation*} \varphi_n(x) = \begin{cases} 0 & \text{if } x \le 0, \\ nx & \text{if } 0 \le x \le \frac1n, \\ 1 & \text{if } x \ge \frac1n. \end{cases} \end{equation*}

Each \(\varphi_n\) is continuous and \(0 \le \varphi_n \le 1\), so \(\varphi_n \in V\).

Claim 1: for each \(k\), the sequence \(\varphi_1(r_k), \varphi_2(r_k), \ldots\) converges, with limit

\begin{equation*} c_k = \begin{cases} 0 & \text{if } r_k \le 0, \\ 1 & \text{if } r_k > 0. \end{cases} \end{equation*}

Indeed, if \(r_k \le 0\) then \(\varphi_n(r_k) = 0\) for every \(n\); and if \(r_k > 0\) then \(\varphi_n(r_k) = 1\) for every \(n > 1/r_k\).

Claim 2: \(\varphi_1, \varphi_2, \ldots\) is a Cauchy sequence in \(V\) with respect to \(\|\cdot\|\). Let \(\varepsilon > 0\). Choose \(K \in \mathbf{Z}^+\) with \(\sum_{k > K} 2^{-k} < \varepsilon^2/2\). Since \(|\varphi_n - \varphi_m| \le 1\) everywhere, the tail of the series contributes at most \(\varepsilon^2/2\) for all \(m, n\). By Claim 1 the finitely many sequences \((\varphi_n(r_k))_n\) for \(k = 1, \ldots, K\) each converge, hence each is Cauchy, so there is \(N\) such that \(|\varphi_n(r_k) - \varphi_m(r_k)|^2 < \varepsilon^2/2\) for all \(m, n \ge N\) and all \(k \le K\). Then for \(m, n \ge N\),

\begin{equation*} \|\varphi_n - \varphi_m\|^2 = \sum_{k=1}^{K} \frac{|\varphi_n(r_k) - \varphi_m(r_k)|^2}{2^k} + \sum_{k > K} \frac{|\varphi_n(r_k) - \varphi_m(r_k)|^2}{2^k} < \frac{\varepsilon^2}{2}\sum_{k=1}^{K} \frac{1}{2^k} + \frac{\varepsilon^2}{2} \le \varepsilon^2. \end{equation*}

Thus \(\|\varphi_n - \varphi_m\| < \varepsilon\) for \(m, n \ge N\), proving Claim 2.

Claim 3: this Cauchy sequence has no limit in \(V\). Suppose \(f \in V\) and \(\|\varphi_n - f\| \to 0\). By the first observation, \(\varphi_n(r_k) \to f(r_k)\) for each \(k\), so by Claim 1, \(f(r_k) = c_k\) for every \(k\). That is, \(f( r) = 0\) for every rational \(r \le 0\) and \(f( r) = 1\) for every rational \(r > 0\).

Choose rationals \(s_n \downarrow 0\) with \(s_n > 0\) and rationals \(t_n \uparrow 0\) with \(t_n \le 0\) (for instance \(s_n = 1/n\) and \(t_n = -1/n\)). Continuity of \(f\) at \(0\) gives

\begin{equation*} f(0) = \lim_{n \to \infty} f(s_n) = 1 \quad \text{and} \quad f(0) = \lim_{n \to \infty} f(t_n) = 0, \end{equation*}

a contradiction. This proves Claim 3.

Hence \(\|\cdot\|\) is not complete on \(V\), so the inner product space of Exercise 1 in Section 8A is not a Hilbert space.

Problem (8B.3)

Suppose \(V_1, V_2, \ldots\) are Hilbert spaces. Let

\begin{equation*} V = \Bigl\{ (f_1, f_2, \ldots) \in V_1 \times V_2 \times \cdots : \sum_{k=1}^{\infty} \|f_k\|^2 < \infty \Bigr\}. \end{equation*}

Show that the equation

\begin{equation*} \langle (f_1, f_2, \ldots), (g_1, g_2, \ldots)\rangle = \sum_{k=1}^{\infty} \langle f_k, g_k\rangle \end{equation*}

defines an inner product on \(V\) that makes \(V\) a Hilbert space.

[Each of the Hilbert spaces \(V_1, V_2, \ldots\) may have a different inner product, even though the same notation is used for the norm and inner product on all these Hilbert spaces.]

Solution

All the \(V_k\) are vector spaces over the same field \(\mathbf{F}\), and \(V_1 \times V_2 \times \cdots\) is a vector space over \(\mathbf{F}\) with coordinatewise operations.

Step 1: \(V\) is a subspace of \(V_1 \times V_2 \times \cdots\).

Clearly \((0, 0, \ldots) \in V\). If \(f = (f_1, f_2, \ldots) \in V\) and \(\alpha \in \mathbf{F}\), then \(\sum_k \|\alpha f_k\|^2 = |\alpha|^2 \sum_k \|f_k\|^2 < \infty\), so \(\alpha f \in V\). If also \(g = (g_1, g_2, \ldots) \in V\), then for each \(k\) the triangle inequality (8.15) and the inequality \((s+t)^2 \le 2s^2 + 2t^2\) for \(s, t \ge 0\) give

\begin{equation*} \|f_k + g_k\|^2 \le \bigl(\|f_k\| + \|g_k\|\bigr)^2 \le 2\|f_k\|^2 + 2\|g_k\|^2 . \end{equation*}

Summing over \(k\) shows \(\sum_k \|f_k + g_k\|^2 < \infty\), so \(f + g \in V\). Thus \(V\) is a vector space.

Step 2: the series defining the inner product converges absolutely.

Suppose \(f, g \in V\). By the Cauchy-Schwarz inequality (8.11) applied in \(V_k\), and again by \(st \le \frac12(s^2 + t^2)\),

\begin{equation*} \bigl|\langle f_k, g_k\rangle\bigr| \le \|f_k\|\,\|g_k\| \le \tfrac12\|f_k\|^2 + \tfrac12\|g_k\|^2 . \end{equation*}

Summing over \(k\), the series \(\sum_{k=1}^{\infty} \langle f_k, g_k \rangle\) converges absolutely, hence converges (in \(\mathbf{F}\), which is complete). So the displayed formula defines a function from \(V \times V\) to \(\mathbf{F}\).

Step 3: it is an inner product.

Positivity: \(\langle f, f\rangle = \sum_{k=1}^{\infty} \langle f_k, f_k\rangle = \sum_{k=1}^{\infty} \|f_k\|^2 \ge 0\).

Definiteness: if \(\langle f, f\rangle = 0\) then \(\sum_k \|f_k\|^2 = 0\), so \(\|f_k\| = 0\) and hence \(f_k = 0\) for every \(k\); thus \(f = 0\). Conversely \(\langle 0, 0\rangle = 0\).

Additivity in the first slot: for \(f, g, h \in V\), all three series \(\sum_k \langle f_k, h_k\rangle\), \(\sum_k \langle g_k, h_k\rangle\), \(\sum_k \langle f_k + g_k, h_k\rangle\) converge by Step 2, and termwise \(\langle f_k + g_k, h_k\rangle = \langle f_k, h_k\rangle + \langle g_k, h_k\rangle\); adding convergent series termwise gives \(\langle f + g, h\rangle = \langle f, h\rangle + \langle g, h\rangle\).

Homogeneity in the first slot: for \(\alpha \in \mathbf{F}\), \(\langle \alpha f, g\rangle = \sum_k \langle \alpha f_k, g_k\rangle = \sum_k \alpha \langle f_k, g_k\rangle = \alpha \langle f, g\rangle\).

Conjugate symmetry: \(\overline{\langle g, f\rangle} = \overline{\sum_k \langle g_k, f_k\rangle} = \sum_k \overline{\langle g_k, f_k\rangle} = \sum_k \langle f_k, g_k\rangle = \langle f, g\rangle\), where conjugation passes through the convergent sum by continuity of conjugation.

Hence \(\langle \cdot, \cdot \rangle\) is an inner product on \(V\), and by 8.19 the associated norm

\begin{equation*} \|f\| = \Bigl(\sum_{k=1}^{\infty} \|f_k\|^2\Bigr)^{1/2} \end{equation*}

is a norm on \(V\). (Here we use the same symbol for the norm on \(V\) and the norms on the \(V_k\); the context makes clear which is meant.)

Step 4: \(V\) is complete.

Suppose \(f^{(1)}, f^{(2)}, \ldots\) is a Cauchy sequence in \(V\); write \(f^{(n)} = (f^{(n)}_1, f^{(n)}_2, \ldots)\).

For each fixed \(k\) and all \(m, n\),

\begin{equation*} \|f^{(n)}_k - f^{(m)}_k\|^2 \le \sum_{j=1}^{\infty} \|f^{(n)}_j - f^{(m)}_j\|^2 = \|f^{(n)} - f^{(m)}\|^2 , \end{equation*}

so \(f^{(1)}_k, f^{(2)}_k, \ldots\) is a Cauchy sequence in \(V_k\). Because \(V_k\) is a Hilbert space, hence complete, there exists \(f_k \in V_k\) with \(\|f^{(n)}_k - f_k\| \to 0\) as \(n \to \infty\). Set \(f = (f_1, f_2, \ldots) \in V_1 \times V_2 \times \cdots\).

Let \(\varepsilon > 0\). Choose \(N\) such that \(\|f^{(n)} - f^{(m)}\| \le \varepsilon\) for all \(m, n \ge N\). Fix \(n \ge N\) and \(K \in \mathbf{Z}^+\). For every \(m \ge N\),

\begin{equation*} \sum_{k=1}^{K} \|f^{(n)}_k - f^{(m)}_k\|^2 \le \|f^{(n)} - f^{(m)}\|^2 \le \varepsilon^2 . \end{equation*}

The left side is a finite sum, and for each \(k\) we have \(f^{(m)}_k \to f_k\) in \(V_k\) as \(m \to \infty\); since the norm is a continuous function (indeed \(\bigl|\,\|u\| - \|v\|\,\bigr| \le \|u - v\|\)), letting \(m \to \infty\) gives

\begin{equation*} \sum_{k=1}^{K} \|f^{(n)}_k - f_k\|^2 \le \varepsilon^2 . \end{equation*}

This holds for every \(K\), so letting \(K \to \infty\),

\begin{equation*} \sum_{k=1}^{\infty} \|f^{(n)}_k - f_k\|^2 \le \varepsilon^2 \qquad \text{for all } n \ge N . \end{equation*}

In particular \(f^{(N)} - f \in V\) (the tuple \((f^{(N)}_1 - f_1, f^{(N)}_2 - f_2, \ldots)\) has square-summable norms), and therefore

\begin{equation*} f = f^{(N)} - \bigl(f^{(N)} - f\bigr) \in V \end{equation*}

because \(V\) is a vector space (Step 1). The displayed inequality now reads \(\|f^{(n)} - f\| \le \varepsilon\) for all \(n \ge N\). As \(\varepsilon > 0\) was arbitrary, \(f^{(n)} \to f\) in \(V\).

Thus every Cauchy sequence in \(V\) converges in \(V\), so \(V\) is a Banach space with the norm coming from the inner product; by 8.21, \(V\) is a Hilbert space.

Problem (8B.4)

Suppose \(V\) is a real Hilbert space. The complexification of \(V\) is the complex vector space \(V_{\mathbf{C}}\) defined by \(V_{\mathbf{C}} = V \times V\), but we write a typical element of \(V_{\mathbf{C}}\) as \(f + ig\) instead of \((f, g)\). Addition and scalar multiplication are defined on \(V_{\mathbf{C}}\) by

\begin{equation*} (f_1 + ig_1) + (f_2 + ig_2) = (f_1 + f_2) + i(g_1 + g_2) \end{equation*}

and

\begin{equation*} (\alpha + i\beta)(f + ig) = (\alpha f - \beta g) + i(\alpha g + \beta f) \end{equation*}

for \(f_1, f_2, f, g_1, g_2, g \in V\) and \(\alpha, \beta \in \mathbf{R}\). Show that

\begin{equation*} \langle f_1 + ig_1, f_2 + ig_2\rangle = \langle f_1, f_2\rangle + \langle g_1, g_2\rangle + i\bigl(\langle g_1, f_2\rangle - \langle f_1, g_2\rangle\bigr) \end{equation*}

defines an inner product on \(V_{\mathbf{C}}\) that makes \(V_{\mathbf{C}}\) into a complex Hilbert space.

Solution

Throughout, \(\langle \cdot, \cdot\rangle\) on the right-hand side denotes the inner product of the real Hilbert space \(V\); it is real valued, symmetric (\(\langle u, v\rangle = \langle v, u\rangle\)) and bilinear over \(\mathbf{R}\). Write \(\|u\| = \langle u, u\rangle^{1/2}\) for \(u \in V\), and denote by \([\cdot, \cdot]\) the proposed form on \(V_{\mathbf{C}}\), so that

\begin{equation*} [f_1 + ig_1, f_2 + ig_2] = \langle f_1, f_2\rangle + \langle g_1, g_2\rangle + i\bigl(\langle g_1, f_2\rangle - \langle f_1, g_2\rangle\bigr) . \end{equation*}

Step 0: \(V_{\mathbf{C}}\) is a complex vector space.

Addition is the coordinatewise addition of \(V \times V\), so \((V_{\mathbf{C}}, +)\) is an abelian group with identity \(0 + i0\). The scalar multiplication is exactly the one obtained by regarding \(V \times V\) as \(\mathbf{R}^2 \otimes V\) with \(\mathbf{C} = \mathbf{R}^2\); concretely, writing \(\lambda = \alpha + i\beta\) and \(\mu = \gamma + i\delta\), we have \(\lambda\mu = (\alpha\gamma - \beta\delta) + i(\alpha\delta + \beta\gamma)\), and

\begin{equation*} \lambda\bigl(\mu(f + ig)\bigr) = \lambda\bigl((\gamma f - \delta g) + i(\gamma g + \delta f)\bigr) = \bigl(\alpha\gamma f - \alpha\delta g - \beta\gamma g - \beta\delta f\bigr) + i\bigl(\alpha\gamma g + \alpha\delta f + \beta\gamma f - \beta\delta g\bigr), \end{equation*}

which equals \((\lambda\mu)(f + ig) = \bigl((\alpha\gamma - \beta\delta)f - (\alpha\delta + \beta\gamma)g\bigr) + i\bigl((\alpha\gamma - \beta\delta)g + (\alpha\delta + \beta\gamma)f\bigr)\). Also \(1(f + ig) = f + ig\), and the two distributive laws follow from the \(\mathbf{R}\)-bilinearity of the formulas. So \(V_{\mathbf{C}}\) is a complex vector space.

Step 1: positivity and definiteness.

Taking \(f_1 = f_2 = f\) and \(g_1 = g_2 = g\) and using the symmetry \(\langle g, f\rangle = \langle f, g\rangle\) of the real inner product,

\begin{equation*} [f + ig, f + ig] = \langle f, f\rangle + \langle g, g\rangle + i\bigl(\langle g, f\rangle - \langle f, g\rangle\bigr) = \|f\|^2 + \|g\|^2 . \end{equation*}

So \([f + ig, f + ig]\) is a nonnegative real number, and it equals \(0\) if and only if \(\|f\| = \|g\| = 0\), that is, if and only if \(f = g = 0\), that is, if and only if \(f + ig\) is the zero element of \(V_{\mathbf{C}}\).

Step 2: additivity in the first slot.

Fix \(f_2 + ig_2 \in V_{\mathbf{C}}\). Each of the four terms \(\langle \cdot, f_2\rangle\), \(\langle \cdot, g_2\rangle\) appearing in the formula is additive in its first argument, and by definition

\begin{equation*} (f_1 + ig_1) + (f_1^{\prime} + ig_1^{\prime}) = (f_1 + f_1^{\prime}) + i(g_1 + g_1^{\prime}) . \end{equation*}

Hence

\begin{equation*} \begin{aligned} \bigl[(f_1 + f_1^{\prime}) + i(g_1 + g_1^{\prime}),\, f_2 + ig_2\bigr] &= \langle f_1 + f_1^{\prime}, f_2\rangle + \langle g_1 + g_1^{\prime}, g_2\rangle + i\bigl(\langle g_1 + g_1^{\prime}, f_2\rangle - \langle f_1 + f_1^{\prime}, g_2\rangle\bigr) \\ &= [f_1 + ig_1, f_2 + ig_2] + [f_1^{\prime} + ig_1^{\prime}, f_2 + ig_2] . \end{aligned} \end{equation*}

Step 3: homogeneity in the first slot for complex scalars.

Let \(\lambda = \alpha + i\beta\) with \(\alpha, \beta \in \mathbf{R}\). Then \(\lambda(f_1 + ig_1) = (\alpha f_1 - \beta g_1) + i(\alpha g_1 + \beta f_1)\), so by definition of \([\cdot,\cdot]\) and \(\mathbf{R}\)-bilinearity of \(\langle \cdot, \cdot\rangle\),

\begin{equation*} \begin{aligned} \bigl[\lambda(f_1 + ig_1),\, f_2 + ig_2\bigr] &= \langle \alpha f_1 - \beta g_1, f_2\rangle + \langle \alpha g_1 + \beta f_1, g_2\rangle \\ &\qquad + i\bigl(\langle \alpha g_1 + \beta f_1, f_2\rangle - \langle \alpha f_1 - \beta g_1, g_2\rangle\bigr) \\ &= \bigl(\alpha\langle f_1, f_2\rangle + \alpha\langle g_1, g_2\rangle - \beta\langle g_1, f_2\rangle + \beta\langle f_1, g_2\rangle\bigr) \\ &\qquad + i\bigl(\alpha\langle g_1, f_2\rangle - \alpha\langle f_1, g_2\rangle + \beta\langle f_1, f_2\rangle + \beta\langle g_1, g_2\rangle\bigr) . \end{aligned} \end{equation*}

On the other hand, setting \(A = \langle f_1, f_2\rangle + \langle g_1, g_2\rangle\) and \(B = \langle g_1, f_2\rangle - \langle f_1, g_2\rangle\) (both real numbers), we have \([f_1 + ig_1, f_2 + ig_2] = A + iB\) and

\begin{equation*} \lambda(A + iB) = (\alpha + i\beta)(A + iB) = (\alpha A - \beta B) + i(\alpha B + \beta A) . \end{equation*}

Expanding, \(\alpha A - \beta B = \alpha\langle f_1, f_2\rangle + \alpha\langle g_1, g_2\rangle - \beta\langle g_1, f_2\rangle + \beta\langle f_1, g_2\rangle\) and \(\alpha B + \beta A = \alpha\langle g_1, f_2\rangle - \alpha\langle f_1, g_2\rangle + \beta\langle f_1, f_2\rangle + \beta\langle g_1, g_2\rangle\). These agree with the two displayed real and imaginary parts above, so

\begin{equation*} \bigl[\lambda(f_1 + ig_1),\, f_2 + ig_2\bigr] = \lambda\,[f_1 + ig_1, f_2 + ig_2] . \end{equation*}

Step 4: conjugate symmetry.

Using symmetry of the real inner product,

\begin{equation*} \begin{aligned} [f_2 + ig_2, f_1 + ig_1] &= \langle f_2, f_1\rangle + \langle g_2, g_1\rangle + i\bigl(\langle g_2, f_1\rangle - \langle f_2, g_1\rangle\bigr) \\ &= \langle f_1, f_2\rangle + \langle g_1, g_2\rangle + i\bigl(\langle f_1, g_2\rangle - \langle g_1, f_2\rangle\bigr) \\ &= A - iB = \overline{[f_1 + ig_1, f_2 + ig_2]} . \end{aligned} \end{equation*}

Steps 1 through 4 show that \([\cdot, \cdot]\) is an inner product on the complex vector space \(V_{\mathbf{C}}\). By 8.19 the associated norm is a norm, and by Step 1 it is given by

\begin{equation*} \|f + ig\|_{\mathbf{C}} = \bigl(\|f\|^2 + \|g\|^2\bigr)^{1/2} . \end{equation*}

Step 5: completeness.

Note first that \(-(f_2 + ig_2) = (-1)(f_2 + ig_2) = (-f_2) + i(-g_2)\), so

\begin{equation*} (f_1 + ig_1) - (f_2 + ig_2) = (f_1 - f_2) + i(g_1 - g_2), \end{equation*}

and therefore

\begin{equation*} \|(f_1 + ig_1) - (f_2 + ig_2)\|_{\mathbf{C}}^2 = \|f_1 - f_2\|^2 + \|g_1 - g_2\|^2 . \end{equation*}

Now suppose \(f_1 + ig_1, f_2 + ig_2, \ldots\) is a Cauchy sequence in \(V_{\mathbf{C}}\). By the displayed identity,

\begin{equation*} \|f_n - f_m\| \le \|(f_n + ig_n) - (f_m + ig_m)\|_{\mathbf{C}} \quad \text{and} \quad \|g_n - g_m\| \le \|(f_n + ig_n) - (f_m + ig_m)\|_{\mathbf{C}} , \end{equation*}

so \(f_1, f_2, \ldots\) and \(g_1, g_2, \ldots\) are Cauchy sequences in \(V\). Since \(V\) is a real Hilbert space, hence complete, there exist \(f, g \in V\) with \(\|f_n - f\| \to 0\) and \(\|g_n - g\| \to 0\). Then

\begin{equation*} \|(f_n + ig_n) - (f + ig)\|_{\mathbf{C}}^2 = \|f_n - f\|^2 + \|g_n - g\|^2 \to 0 . \end{equation*}

So the sequence converges in \(V_{\mathbf{C}}\). Hence \(V_{\mathbf{C}}\) is a complex Banach space with the norm coming from \([\cdot, \cdot]\), and by 8.21 it is a complex Hilbert space.

Problem (8B.5)

Prove that if \(V\) is a normed vector space, \(f \in V\), and \(r > 0\), then the open ball \(B(f, r)\) centered at \(f\) with radius \(r\) is convex.

Solution

Recall \(B(f, r) = \{h \in V : \|h - f\| < r\}\), and recall from 8.25 that a subset \(U\) of a vector space is convex if \((1-t)g + th \in U\) for all \(t \in [0,1]\) and all \(g, h \in U\).

Suppose \(g, h \in B(f, r)\) and \(t \in [0,1]\). Because \((1-t) + t = 1\) we may write

\begin{equation*} \bigl((1-t)g + th\bigr) - f = (1-t)(g - f) + t(h - f) . \end{equation*}

The triangle inequality and homogeneity of the norm now give

\begin{equation*} \bigl\|\bigl((1-t)g + th\bigr) - f\bigr\| \le (1-t)\|g - f\| + t\,\|h - f\| , \end{equation*}

where we used \(|1-t| = 1-t\) and \(|t| = t\) since \(t \in [0,1]\).

We claim the right side is strictly less than \(r\). Since \(\|g - f\| < r\) and \(\|h - f\| < r\), we have \((1-t)\|g-f\| \le (1-t)r\) with strict inequality unless \(t = 1\), and \(t\|h - f\| \le tr\) with strict inequality unless \(t = 0\). Because \(t\) cannot equal both \(0\) and \(1\), at least one of these two inequalities is strict, and hence

\begin{equation*} (1-t)\|g - f\| + t\,\|h - f\| < (1-t)r + tr = r . \end{equation*}

Combining the last two displays gives \(\bigl\|\bigl((1-t)g + th\bigr) - f\bigr\| < r\), so \((1-t)g + th \in B(f, r)\). Since \(g, h \in B(f,r)\) and \(t \in [0,1]\) were arbitrary, \(B(f, r)\) is convex.

Problem (8B.6)

(a) Suppose \(V\) is an inner product space and \(B\) is the open unit ball in \(V\) (thus \(B = \{f \in V : \|f\| < 1\}\)). Prove that if \(U\) is a subset of \(V\) such that \(B \subseteq U \subseteq \overline{B}\), then \(U\) is convex.

(b) Give an example to show that (a) can fail if the phrase inner product space is replaced by Banach space.

Solution

(a) First we identify \(\overline{B}\). In any normed vector space, \(\overline{B} = \{f \in V : \|f\| \le 1\}\). Indeed, the set \(\{f : \|f\| \le 1\}\) is closed (it is the inverse image of the closed set \([0,1]\) under the continuous function \(\|\cdot\|\)) and contains \(B\), so it contains \(\overline{B}\). Conversely if \(\|f\| \le 1\), then the elements \((1 - \frac1n)f\) lie in \(B\) (their norms are at most \(1 - \frac1n < 1\)) and converge to \(f\), so \(f \in \overline{B}\).

Now suppose \(B \subseteq U \subseteq \overline{B}\), let \(f, g \in U\), let \(t \in [0,1]\), and set

\begin{equation*} h = (1-t)f + tg . \end{equation*}

We must show \(h \in U\). If \(t = 0\) then \(h = f \in U\), and if \(t = 1\) then \(h = g \in U\); so assume \(0 < t < 1\).

Since \(f, g \in U \subseteq \overline{B}\) we have \(\|f\| \le 1\) and \(\|g\| \le 1\), so by the triangle inequality (8.15)

\begin{equation*} \|h\| \le (1-t)\|f\| + t\|g\| \le (1-t) + t = 1 . \end{equation*}

Case 1: \(\|h\| < 1\). Then \(h \in B \subseteq U\), as desired.

Case 2: \(\|h\| = 1\). Then equality holds throughout the display above, so in particular

\begin{equation*} (1-t)\|f\| + t\|g\| = 1 . \end{equation*}

Since \(\|f\| \le 1\), \(\|g\| \le 1\), \(1 - t > 0\) and \(t > 0\), this forces \(\|f\| = \|g\| = 1\); otherwise the left side would be strictly less than \((1-t) + t = 1\).

Set \(u = (1-t)f\) and \(v = tg\), so \(h = u + v\), \(\|u\| = 1 - t > 0\), \(\|v\| = t > 0\), and

\begin{equation*} \|u + v\| = \|h\| = 1 = (1-t) + t = \|u\| + \|v\| . \end{equation*}

Thus equality holds in the triangle inequality for \(u\) and \(v\). As shown in the discussion following 8.15 (see 8.18), equality in the triangle inequality holds if and only if \(\langle u, v\rangle = \|u\|\,\|v\|\), and this happens if and only if one of \(u, v\) is a nonnegative multiple of the other. Since \(u \ne 0\) and \(v \ne 0\), in either case there is \(c > 0\) with \(v = cu\). Taking norms, \(t = \|v\| = c\|u\| = c(1-t)\), so \(c = t/(1-t)\) and

\begin{equation*} tg = v = cu = \frac{t}{1-t}\,(1-t)f = tf . \end{equation*}

Since \(t \ne 0\) this gives \(g = f\), and therefore \(h = (1-t)f + tf = f \in U\).

In both cases \(h \in U\), so \(U\) is convex.

(b) Let \(V = \mathbf{R}^2\) with the norm

\begin{equation*} \|(x_1, x_2)\| = \max\{|x_1|, |x_2|\} . \end{equation*}

This is a norm, and \(V\) is a Banach space: a sequence in \(V\) is Cauchy (respectively convergent) with respect to this norm exactly when each of the two coordinate sequences is Cauchy (respectively convergent) in \(\mathbf{R}\), because \(|x_j| \le \|(x_1, x_2)\| \le |x_1| + |x_2|\); completeness of \(\mathbf{R}\) then gives completeness of \(V\).

Here

\begin{equation*} B = \{(x_1, x_2) : |x_1| < 1 \text{ and } |x_2| < 1\}, \qquad \overline{B} = \{(x_1, x_2) : |x_1| \le 1 \text{ and } |x_2| \le 1\} , \end{equation*}

the open and closed unit squares. Let

\begin{equation*} U = B \cup \{(1, 1), (1, -1)\} . \end{equation*}

Then \(B \subseteq U \subseteq \overline{B}\), since \(\|(1,1)\| = \|(1,-1)\| = 1\).

However \(U\) is not convex: the points \((1,1)\) and \((1,-1)\) belong to \(U\), but

\begin{equation*} \tfrac12 (1,1) + \tfrac12 (1,-1) = (1, 0) , \end{equation*}

and \((1,0) \notin B\) because \(\|(1,0)\| = 1\), while \((1,0)\) is neither \((1,1)\) nor \((1,-1)\). Hence \((1,0) \notin U\), so \(U\) fails to be convex. This shows that (a) can fail when the inner product space hypothesis is weakened to a Banach space.

(The mechanism is exactly what the proof of (a) used: in an inner product space the unit sphere contains no line segment, whereas the unit sphere of this Banach space contains the segment from \((1,-1)\) to \((1,1)\).)

Problem (8B.7)

Suppose \(V\) is a normed vector space and \(U\) is a closed subset of \(V\). Prove that \(U\) is convex if and only if

\begin{equation*} \frac{f + g}{2} \in U \quad \text{for all } f, g \in U . \end{equation*}

Solution

One direction is immediate. If \(U\) is convex and \(f, g \in U\), then taking \(t = \frac12\) in the definition of convexity (8.25) gives \(\frac12 f + \frac12 g = \frac{f+g}{2} \in U\). (This direction does not use that \(U\) is closed.)

For the other direction, suppose \(\frac{f+g}{2} \in U\) for all \(f, g \in U\). Fix \(f, g \in U\) and define

\begin{equation*} D = \bigl\{ t \in [0,1] : (1-t)f + tg \in U \bigr\} . \end{equation*}

We must show \(D = [0,1]\).

Step 1: \(D\) is closed under averaging. Suppose \(s, t \in D\). Then \((1-s)f + sg \in U\) and \((1-t)f + tg \in U\), so by hypothesis their average lies in \(U\). That average is

\begin{equation*} \frac{\bigl((1-s)f + sg\bigr) + \bigl((1-t)f + tg\bigr)}{2} = \Bigl(1 - \frac{s+t}{2}\Bigr) f + \frac{s+t}{2}\, g , \end{equation*}

and \(\frac{s+t}{2} \in [0,1]\). Hence \(\frac{s+t}{2} \in D\).

Step 2: \(D\) contains every dyadic rational in \([0,1]\). Call a number of the form \(k/2^n\) with \(n \ge 0\) an integer and \(0 \le k \le 2^n\) a dyadic rational in \([0,1]\); we induct on \(n\).

For \(n = 0\) the only such numbers are \(0\) and \(1\), and \(0 \in D\) (since \((1-0)f + 0g = f \in U\)) and \(1 \in D\) (since \(g \in U\)).

Suppose \(n \ge 0\) and \(k/2^n \in D\) for all integers \(k\) with \(0 \le k \le 2^n\). Let \(j\) be an integer with \(0 \le j \le 2^{n+1}\). If \(j\) is even, say \(j = 2k\), then \(j/2^{n+1} = k/2^n \in D\) by the induction hypothesis. If \(j\) is odd, say \(j = 2k+1\) with \(0 \le k \le 2^n - 1\), then

\begin{equation*} \frac{j}{2^{n+1}} = \frac{1}{2}\Bigl(\frac{k}{2^n} + \frac{k+1}{2^n}\Bigr) , \end{equation*}

and both \(k/2^n\) and \((k+1)/2^n\) lie in \(D\) by the induction hypothesis, so Step 1 gives \(j/2^{n+1} \in D\). This completes the induction.

Step 3: \(D\) is a closed subset of \([0,1]\). Define \(\gamma : [0,1] \to V\) by \(\gamma(t) = (1-t)f + tg\). For \(s, t \in [0,1]\),

\begin{equation*} \|\gamma(t) - \gamma(s)\| = \|(t-s)(g - f)\| = |t - s|\,\|g - f\| , \end{equation*}

so \(\gamma\) is continuous. If \(t_1, t_2, \ldots\) is a sequence in \(D\) converging to \(t \in [0,1]\), then \(\gamma(t_n) \to \gamma(t)\) in \(V\), and each \(\gamma(t_n) \in U\); because \(U\) is closed, \(\gamma(t) \in U\), that is, \(t \in D\). Hence \(D\) is closed in \([0,1]\).

Step 4: conclusion. The dyadic rationals are dense in \([0,1]\): given \(t \in [0,1]\) and \(n \ge 0\), the number \(\lfloor 2^n t\rfloor / 2^n\) is a dyadic rational in \([0,1]\) within \(2^{-n}\) of \(t\). By Step 2 every dyadic rational in \([0,1]\) lies in \(D\), so every \(t \in [0,1]\) is a limit of a sequence of points of \(D\); by Step 3, \(t \in D\). Therefore \(D = [0,1]\), which says precisely that

\begin{equation*} (1-t)f + tg \in U \quad \text{for all } t \in [0,1] . \end{equation*}

Since \(f, g \in U\) were arbitrary, \(U\) is convex.

Problem (8B.8)

Prove that if \(U\) is a convex subset of a normed vector space, then \(\overline{U}\) is also convex.

Solution

Suppose \(V\) is a normed vector space and \(U \subseteq V\) is convex. To show that \(\overline{U}\) is convex we must show, following 8.25, that

\begin{equation*} (1-t) f + t g \in \overline{U} \quad \text{for all } t \in [0,1] \text{ and all } f, g \in \overline{U}. \end{equation*}

So suppose \(f, g \in \overline{U}\) and \(t \in [0,1]\). Because \(\overline{U}\) is the closure of \(U\), there exist sequences \(f_1, f_2, \ldots\) and \(g_1, g_2, \ldots\) in \(U\) with

\begin{equation*} \lim_{k \to \infty} \|f - f_k\| = 0 \quad \text{and} \quad \lim_{k \to \infty} \|g - g_k\| = 0 . \end{equation*}

(Indeed, for each \(k \in \mathbf{Z}^+\) the ball \(B(f, 1/k)\) intersects \(U\), because otherwise \(f\) would lie in the open set \(V \setminus \overline{U}\); choose \(f_k\) in that intersection, and similarly for \(g_k\).)

For each \(k \in \mathbf{Z}^+\) set

\begin{equation*} h_k = (1-t) f_k + t g_k . \end{equation*}

Because \(f_k, g_k \in U\) and \(t \in [0,1]\) and \(U\) is convex, we have \(h_k \in U\) for each \(k\).

Let \(h = (1-t) f + t g\). Then the triangle inequality and the homogeneity of the norm give

\begin{equation*} \|h - h_k\| = \|(1-t)(f - f_k) + t (g - g_k)\| \le (1-t)\,\|f - f_k\| + t\,\|g - g_k\| \end{equation*}

for each \(k \in \mathbf{Z}^+\), where we have used \(1 - t \ge 0\) and \(t \ge 0\). The right side of the inequality above has limit \(0\) as \(k \to \infty\). Hence \(\lim_{k \to \infty} \|h - h_k\| = 0\).

Thus \(h\) is the limit of a sequence of elements of \(U\), and therefore \(h \in \overline{U}\). In other words, \((1-t) f + t g \in \overline{U}\). Because \(f, g \in \overline{U}\) and \(t \in [0,1]\) were arbitrary, \(\overline{U}\) is convex.

Problem (8B.9)

Prove that if \(U\) is a convex subset of a normed vector space, then the interior of \(U\) is also convex.

[The interior of \(U\) is the set \(\{f \in U : B(f, r) \subseteq U \text{ for some } r > 0\}\).]

Solution

Suppose \(V\) is a normed vector space and \(U \subseteq V\) is convex. Write

\begin{equation*} \operatorname{int} U = \{f \in U : B(f, r) \subseteq U \text{ for some } r > 0\}, \end{equation*}

where \(B(f, r) = \{k \in V : \|k - f\| < r\}\) is the open ball of radius \(r\) centered at \(f\).

If \(\operatorname{int} U\) is empty then it is convex (the condition in 8.25 holds vacuously). So suppose \(f, g \in \operatorname{int} U\) and \(t \in [0,1]\); we must show that

\begin{equation*} h = (1-t) f + t g \end{equation*}

belongs to \(\operatorname{int} U\).

Choose \(r > 0\) and \(s > 0\) such that \(B(f, r) \subseteq U\) and \(B(g, s) \subseteq U\), and set \(\rho = \min\{r, s\} > 0\). We claim that

\begin{equation*} B(h, \rho) \subseteq U , \end{equation*}

which will show that \(h \in \operatorname{int} U\) (note that \(h \in B(h,\rho) \subseteq U\), so in particular \(h \in U\)).

To prove the claim, suppose \(k \in B(h, \rho)\) and put \(u = k - h\), so \(\|u\| < \rho\). Then

\begin{equation*} \|(f + u) - f\| = \|u\| < \rho \le r \quad \text{and} \quad \|(g + u) - g\| = \|u\| < \rho \le s , \end{equation*}

so \(f + u \in B(f, r) \subseteq U\) and \(g + u \in B(g, s) \subseteq U\). Because \(U\) is convex and \(t \in [0,1]\), we conclude that

\begin{equation*} (1-t)(f + u) + t (g + u) \in U . \end{equation*}

However

\begin{equation*} (1-t)(f + u) + t (g + u) = \bigl((1-t) f + t g\bigr) + \bigl((1-t) + t\bigr) u = h + u = k . \end{equation*}

Hence \(k \in U\), proving the claim.

Thus \(h \in \operatorname{int} U\). Because \(f, g \in \operatorname{int} U\) and \(t \in [0,1]\) were arbitrary, \(\operatorname{int} U\) is convex.

Problem (8B.10)

Suppose \(V\) is a Hilbert space, \(U\) is a nonempty closed convex subset of \(V\), and \(g \in U\) is the unique element of \(U\) with smallest norm (obtained by taking \(f = 0\) in 8.28). Prove that

\begin{equation*} \operatorname{Re} \langle g, h \rangle \ge \|g\|^2 \end{equation*}

for all \(h \in U\).

Solution

Taking \(f = 0\) in 8.28 shows that there is indeed a unique \(g \in U\) with

\begin{equation*} \|g\| = \|0 - g\| = \operatorname{distance}(0, U) = \inf\{\|k\| : k \in U\} ; \end{equation*}

that is, \(g\) is the unique element of \(U\) of smallest norm. In particular \(\|g\| \le \|k\|\) for every \(k \in U\); call this the minimality property of \(g\).

Now suppose \(h \in U\). Let \(t \in (0, 1]\). Because \(g, h \in U\) and \(U\) is convex, 8.25 gives

\begin{equation*} (1 - t) g + t h \in U . \end{equation*}

Note that \((1 - t) g + t h = g + t (h - g)\). Hence the minimality property of \(g\), applied to this element of \(U\), gives

\begin{equation*} \|g\|^2 \le \|g + t(h - g)\|^2 . \end{equation*}

Expanding the right side using the properties of the inner product (and using that \(t\) is real) gives

\begin{equation*} \begin{aligned} \|g + t(h-g)\|^2 &= \langle g + t(h-g),\, g + t(h-g) \rangle \\ &= \|g\|^2 + t \langle h - g, g \rangle + t \langle g, h - g \rangle + t^2 \|h - g\|^2 \\ &= \|g\|^2 + 2 t \operatorname{Re}\langle g, h - g \rangle + t^2 \|h - g\|^2 , \end{aligned} \end{equation*}

where the last line holds because \(\langle h - g, g\rangle = \overline{\langle g, h - g\rangle}\) and \(z + \bar z = 2\operatorname{Re} z\) for every \(z \in \mathbf{C}\).

Combining the last two displays and subtracting \(\|g\|^2\) from both sides, we obtain

\begin{equation*} 0 \le 2 t \operatorname{Re}\langle g, h - g \rangle + t^2 \|h - g\|^2 . \end{equation*}

Because \(t > 0\), we may divide by \(2t\), getting

\begin{equation*} 0 \le \operatorname{Re}\langle g, h - g \rangle + \tfrac{t}{2} \|h - g\|^2 \end{equation*}

for every \(t \in (0, 1]\). Letting \(t \to 0^+\) (the left-hand term does not depend on \(t\)) gives

\begin{equation*} 0 \le \operatorname{Re}\langle g, h - g \rangle . \end{equation*}

Finally, additivity of the inner product in its first slot and the identity \(\langle g, g\rangle = \|g\|^2\) give

\begin{equation*} \operatorname{Re}\langle g, h - g\rangle = \operatorname{Re}\langle g, h \rangle - \operatorname{Re}\langle g, g \rangle = \operatorname{Re}\langle g, h \rangle - \|g\|^2 . \end{equation*}

Hence \(\operatorname{Re}\langle g, h \rangle - \|g\|^2 \ge 0\), which is the desired inequality

\begin{equation*} \operatorname{Re}\langle g, h \rangle \ge \|g\|^2 . \end{equation*}

Because \(h \in U\) was arbitrary, the proof is complete.

[Geometrically: the closed half-space \(\{k \in V : \operatorname{Re}\langle g, k\rangle \ge \|g\|^2\}\) contains \(U\) and touches \(U\) at \(g\); its boundary hyperplane is perpendicular to \(g\) and supports \(U\) at the point of \(U\) closest to the origin.]

Problem (8B.11)

Suppose \(V\) is a Hilbert space. A closed half-space of \(V\) is a set of the form

\begin{equation*} \{g \in V : \operatorname{Re}\langle g, h \rangle \ge c\} \end{equation*}

for some \(h \in V\) and some \(c \in \mathbf{R}\). Prove that every closed convex subset of \(V\) is the intersection of all the closed half-spaces that contain it.

Solution

Suppose \(C \subseteq V\) is closed and convex. Let \(\mathcal{H}\) denote the collection of all closed half-spaces of \(V\) that contain \(C\), and let

\begin{equation*} I = \bigcap_{H \in \mathcal{H}} H \end{equation*}

(with the usual convention that an intersection over the empty collection equals \(V\)). We must prove that \(C = I\).

The inclusion \(C \subseteq I\) is immediate: every \(H \in \mathcal{H}\) contains \(C\) by definition of \(\mathcal{H}\), so \(C\) is contained in their intersection. The work is in proving \(I \subseteq C\), which we do by showing that if \(f \in V\) and \(f \notin C\), then some closed half-space contains \(C\) but not \(f\).

The case \(C = \varnothing\). Take \(h = 0\) and \(c = 1\) in the definition of closed half-space. Then

\begin{equation*} \{g \in V : \operatorname{Re}\langle g, 0 \rangle \ge 1\} = \{g \in V : 0 \ge 1\} = \varnothing , \end{equation*}

so \(\varnothing\) is itself a closed half-space, and it contains \(C = \varnothing\). Hence \(\varnothing \in \mathcal{H}\) and therefore \(I = \varnothing = C\), as desired.

The case \(C \ne \varnothing\). Suppose \(f \in V\) with \(f \notin C\). Because \(V\) is a Hilbert space and \(C\) is a nonempty closed convex subset of \(V\), 8.28 provides \(g \in C\) with

\begin{equation*} \|f - g\| = \operatorname{distance}(f, C) . \end{equation*}

Set \(u = f - g\). Because \(g \in C\) and \(f \notin C\), we have \(u \ne 0\), so \(\|u\|^2 > 0\).

We first prove the variational inequality

\begin{equation*} \operatorname{Re}\langle u, k - g \rangle \le 0 \quad \text{for every } k \in C . \end{equation*}

Indeed, let \(k \in C\) and let \(t \in (0,1]\). Since \(g, k \in C\) and \(C\) is convex, 8.25 gives \((1-t) g + t k = g + t(k - g) \in C\). Because \(g\) minimizes the distance from \(f\) to \(C\),

\begin{equation*} \|u\|^2 = \|f - g\|^2 \le \bigl\| f - g - t(k-g) \bigr\|^2 = \|u\|^2 - 2t \operatorname{Re}\langle u, k - g\rangle + t^2 \|k - g\|^2 , \end{equation*}

where the last equality comes from expanding

\begin{equation*} \|u - t(k-g)\|^2 = \langle u - t(k-g),\, u - t(k-g)\rangle = \|u\|^2 - t\langle u, k-g\rangle - t \langle k-g, u\rangle + t^2\|k-g\|^2 \end{equation*}

and then using \(\langle k-g, u\rangle = \overline{\langle u, k-g\rangle}\) together with \(z + \bar z = 2\operatorname{Re} z\) (the scalar \(t\) is real, so it passes through both slots unconjugated, by 8.3(c)). Subtracting \(\|u\|^2\) and dividing by \(2t > 0\) gives

\begin{equation*} \operatorname{Re}\langle u, k - g \rangle \le \tfrac{t}{2} \|k - g\|^2 \end{equation*}

for all \(t \in (0,1]\). Letting \(t \to 0^+\) yields \(\operatorname{Re}\langle u, k - g\rangle \le 0\), as claimed.

Now define

\begin{equation*} h = -u \quad \text{and} \quad c = -\operatorname{Re}\langle g, u \rangle , \end{equation*}

and let \(H = \{k \in V : \operatorname{Re}\langle k, h \rangle \ge c\}\), a closed half-space of \(V\).

\(H\) contains \(C\): if \(k \in C\), then the variational inequality together with \(\operatorname{Re}\langle k, u\rangle = \operatorname{Re}\langle u, k\rangle\) (conjugates have equal real parts) gives

\begin{equation*} \operatorname{Re}\langle k, u \rangle - \operatorname{Re}\langle g, u \rangle = \operatorname{Re}\langle u, k - g \rangle \le 0 , \end{equation*}

hence \(\operatorname{Re}\langle k, u\rangle \le \operatorname{Re}\langle g, u\rangle\), hence

\begin{equation*} \operatorname{Re}\langle k, h \rangle = -\operatorname{Re}\langle k, u \rangle \ge -\operatorname{Re}\langle g, u \rangle = c . \end{equation*}

Thus \(k \in H\), and so \(C \subseteq H\); that is, \(H \in \mathcal{H}\).

\(H\) does not contain \(f\): because \(f = g + u\),

\begin{equation*} \operatorname{Re}\langle f, h \rangle = -\operatorname{Re}\langle g + u, u \rangle = -\operatorname{Re}\langle g, u \rangle - \|u\|^2 = c - \|u\|^2 < c , \end{equation*}

where the strict inequality holds because \(\|u\|^2 > 0\). Hence \(f \notin H\).

Since \(H \in \mathcal{H}\) and \(f \notin H\), we conclude \(f \notin I\). We have therefore shown that \(V \setminus C \subseteq V \setminus I\), which is equivalent to \(I \subseteq C\).

Combining the two inclusions gives \(C = I\): every closed convex subset of \(V\) is the intersection of all the closed half-spaces that contain it.

Problem (8B.12)

Give an example of a nonempty closed subset \(U\) of the Hilbert space \(\ell^2\) and \(a \in \ell^2\) such that there does not exist \(b \in U\) with \(\|a - b\| = \operatorname{distance}(a, U)\).

[By 8.28, \(U\) cannot be a convex subset of \(\ell^2\).]

Solution

For \(k \in \mathbf{Z}^+\) let \(e_k \in \ell^2\) denote the sequence whose \(k^{\text{th}}\) coordinate is \(1\) and whose other coordinates are \(0\). Take

\begin{equation*} a = 0 \quad \text{and} \quad U = \Bigl\{ \bigl(1 + \tfrac{1}{k}\bigr) e_k : k \in \mathbf{Z}^+ \Bigr\} . \end{equation*}

Write \(u_k = \bigl(1 + \frac{1}{k}\bigr) e_k\), so \(U = \{u_1, u_2, \ldots\}\). Clearly \(U\) is a nonempty subset of \(\ell^2\) (each \(u_k\) has only one nonzero coordinate, so \(\|u_k\| = 1 + \frac1k < \infty\)).

\(U\) is closed. If \(j \ne k\), then \(u_j\) and \(u_k\) are supported on disjoint coordinates, so

\begin{equation*} \|u_j - u_k\|^2 = \Bigl(1 + \tfrac{1}{j}\Bigr)^2 + \Bigl(1 + \tfrac{1}{k}\Bigr)^2 > 1 + 1 = 2 , \end{equation*}

and hence \(\|u_j - u_k\| > \sqrt{2}\). Now suppose \(f_1, f_2, \ldots\) is a sequence in \(U\) converging to some \(f \in \ell^2\). A convergent sequence is Cauchy, so there exists \(N \in \mathbf{Z}^+\) such that \(\|f_m - f_n\| < \sqrt{2}\) for all \(m, n \ge N\). By the displayed estimate this forces \(f_m = f_n\) for all \(m, n \ge N\); that is, the sequence is eventually equal to the constant \(f_N \in U\). Hence \(f = f_N \in U\). Thus \(U\) contains all limits of convergent sequences of its elements, so \(U\) is closed.

The distance is not attained. We have

\begin{equation*} \operatorname{distance}(a, U) = \inf_{k \in \mathbf{Z}^+} \|0 - u_k\| = \inf_{k \in \mathbf{Z}^+} \Bigl(1 + \tfrac{1}{k}\Bigr) = 1 . \end{equation*}

However, for each \(k \in \mathbf{Z}^+\) we have \(\|a - u_k\| = 1 + \frac{1}{k} > 1\). Since every element of \(U\) is of the form \(u_k\), there is no \(b \in U\) with \(\|a - b\| = \operatorname{distance}(a, U)\).

As the hint notes, this is consistent with 8.28 only because \(U\) is not convex: indeed \(U\) is a discrete set of isolated points, and for instance \(\frac{1}{2} u_1 + \frac{1}{2} u_2 = (1, \frac{3}{4}, 0, 0, \ldots) \notin U\).

Problem (8B.13)

In the real Banach space \(\mathbf{R}^2\) with norm defined by \(\|(x, y)\|_\infty = \max\{|x|, |y|\}\), give an example of a closed convex set \(U \subseteq \mathbf{R}^2\) and \(z \in \mathbf{R}^2\) such that there exist infinitely many choices of \(w \in U\) with \(\|z - w\|_\infty = \operatorname{distance}(z, U)\).

Solution

Take

\begin{equation*} U = \{(x, y) \in \mathbf{R}^2 : x \ge 1\} \quad \text{and} \quad z = (0, 0) . \end{equation*}

\(U\) is closed. \(U\) is the inverse image of the closed set \([1, \infty)\) under the map \((x,y) \mapsto x\), which is continuous with respect to \(\|\cdot\|_\infty\) because \(|x - x^{\prime}| \le \|(x,y) - (x^{\prime},y^{\prime})\|_\infty\). Hence \(U\) is closed.

\(U\) is convex. If \((x, y), (x^{\prime}, y^{\prime}) \in U\) and \(t \in [0,1]\), then the first coordinate of \((1-t)(x,y) + t(x^{\prime},y^{\prime})\) equals \((1-t) x + t x^{\prime} \ge (1-t) + t = 1\), so this point lies in \(U\). Thus \(U\) is convex, as required by 8.25.

The distance from \(z\) to \(U\) equals \(1\). If \(w = (x, y) \in U\), then

\begin{equation*} \|z - w\|_\infty = \max\{|x|, |y|\} \ge |x| = x \ge 1 , \end{equation*}

using \(x \ge 1 > 0\). Hence \(\operatorname{distance}(z, U) \ge 1\). On the other hand \((1, 0) \in U\) and \(\|z - (1,0)\|_\infty = 1\), so \(\operatorname{distance}(z, U) = 1\).

Infinitely many closest points. For every \(y \in \mathbf{R}\) with \(|y| \le 1\), the point \(w = (1, y)\) lies in \(U\) and satisfies

\begin{equation*} \|z - w\|_\infty = \max\{1, |y|\} = 1 = \operatorname{distance}(z, U) . \end{equation*}

There are infinitely many such \(y\) (indeed a whole interval of them), so there are infinitely many \(w \in U\) attaining the distance.

For completeness, the set of closest points is exactly \(\{(1, y) : |y| \le 1\}\): if \(w = (x,y) \in U\) with \(\|z - w\|_\infty = 1\), then \(x \le \max\{|x|,|y|\} = 1\) and \(x \ge 1\) force \(x = 1\), and \(|y| \le \max\{|x|,|y|\} = 1\).

[This shows that the uniqueness assertion of 8.28 genuinely requires an inner product: \(\mathbf{R}^2\) with \(\|\cdot\|_\infty\) is a Banach space, and \(U\) here is even a closed convex set with nonempty interior, yet the nearest point is far from unique. The failure traces back to the parallelogram equality 8.20, which fails for \(\|\cdot\|_\infty\) and which drove the uniqueness argument in the proof of 8.28.]

Problem (8B.14)

Suppose \(f\) and \(g\) are elements of an inner product space. Prove that \(\langle f, g \rangle = 0\) if and only if

\begin{equation*} \|f\| \le \|f + \alpha g\| \end{equation*}

for all \(\alpha \in \mathbf{F}\).

Solution

Throughout, \(V\) denotes the inner product space containing \(f\) and \(g\), and we use the expansion

\begin{equation*} \|f + \alpha g\|^2 = \langle f + \alpha g, f + \alpha g \rangle = \|f\|^2 + 2 \operatorname{Re}\bigl(\overline{\alpha}\,\langle f, g \rangle\bigr) + |\alpha|^2 \|g\|^2 \end{equation*}

valid for all \(\alpha \in \mathbf{F}\). This follows from additivity in each slot together with \(\langle f, \alpha g\rangle = \overline{\alpha} \langle f, g\rangle\) and \(\langle \alpha g, f \rangle = \alpha \langle g, f\rangle = \alpha \overline{\langle f, g\rangle} = \overline{\overline{\alpha}\langle f,g\rangle}\), so that the two cross terms sum to \(2 \operatorname{Re}(\overline{\alpha}\langle f, g\rangle)\). (When \(\mathbf{F} = \mathbf{R}\) the complex conjugations are vacuous and the formula reads \(\|f\|^2 + 2\alpha \langle f, g\rangle + \alpha^2\|g\|^2\).)

First direction. Suppose \(\langle f, g \rangle = 0\). Then for each \(\alpha \in \mathbf{F}\) we have \(\langle f, \alpha g\rangle = \overline{\alpha}\langle f, g\rangle = 0\), so \(f\) and \(\alpha g\) are orthogonal. The Pythagorean Theorem (8.9) then gives

\begin{equation*} \|f + \alpha g\|^2 = \|f\|^2 + \|\alpha g\|^2 = \|f\|^2 + |\alpha|^2 \|g\|^2 \ge \|f\|^2 . \end{equation*}

Taking square roots (norms are nonnegative) gives \(\|f\| \le \|f + \alpha g\|\) for all \(\alpha \in \mathbf{F}\), as desired.

Second direction. Now suppose \(\|f\| \le \|f + \alpha g\|\) for all \(\alpha \in \mathbf{F}\); we must show \(\langle f, g\rangle = 0\).

If \(g = 0\), then \(\langle f, g\rangle = \langle f, 0\rangle = 0\) and we are done. So assume \(g \ne 0\), which gives \(\|g\|^2 > 0\), and set

\begin{equation*} \alpha = -\frac{\langle f, g \rangle}{\|g\|^2} \in \mathbf{F} . \end{equation*}

For this \(\alpha\) we compute

\begin{equation*} \overline{\alpha}\,\langle f, g \rangle = -\frac{\overline{\langle f, g\rangle}\,\langle f, g\rangle}{\|g\|^2} = -\frac{|\langle f, g\rangle|^2}{\|g\|^2} \end{equation*}

(a real number, so it equals its own real part), and

\begin{equation*} |\alpha|^2 \|g\|^2 = \frac{|\langle f, g\rangle|^2}{\|g\|^4}\,\|g\|^2 = \frac{|\langle f, g\rangle|^2}{\|g\|^2} . \end{equation*}

Substituting these into the expansion at the start of this solution gives

\begin{equation*} \|f + \alpha g\|^2 = \|f\|^2 - 2\,\frac{|\langle f, g\rangle|^2}{\|g\|^2} + \frac{|\langle f, g\rangle|^2}{\|g\|^2} = \|f\|^2 - \frac{|\langle f, g\rangle|^2}{\|g\|^2} . \end{equation*}

The hypothesis \(\|f\| \le \|f + \alpha g\|\), squared, now yields

\begin{equation*} \|f\|^2 \le \|f\|^2 - \frac{|\langle f, g\rangle|^2}{\|g\|^2} , \end{equation*}

and hence

\begin{equation*} \frac{|\langle f, g\rangle|^2}{\|g\|^2} \le 0 . \end{equation*}

Because \(\|g\|^2 > 0\) and \(|\langle f, g\rangle|^2 \ge 0\), this forces \(|\langle f, g\rangle|^2 = 0\), so \(\langle f, g \rangle = 0\).

Both directions are proved, so \(\langle f, g\rangle = 0\) if and only if \(\|f\| \le \|f + \alpha g\|\) for all \(\alpha \in \mathbf{F}\).

[This characterization is exactly the mechanism behind 8.37(a): there one knows that \(\|f - P_U f\| \le \|f - P_U f + \alpha g\|\) for all \(g \in U\) and all \(\alpha \in \mathbf{F}\), because \(P_U f + \alpha g\) ranges over elements of the subspace \(U\) and \(P_U f\) is the closest such element; the equivalence just proved then says that \(f - P_U f\) is orthogonal to every \(g \in U\).]

Problem (8B.15)

Suppose \(U\) is a closed subspace of a Hilbert space \(V\) and \(f \in V\). Prove that

\begin{equation*} \|P_U f\| \leq \|f\|, \end{equation*}

with equality if and only if \(f \in U\).

[This exercise asks you to prove 8.37(d).]

Solution

Because \(U\) is a closed subspace of \(V\), it is in particular a nonempty closed convex subset of \(V\), so \(P_U f\) is defined (8.34, which makes sense by 8.28).

By 8.37(a), the element \(f - P_U f\) is orthogonal to every element of \(U\). Since \(P_U f \in U\), this gives

\begin{equation*} \langle f - P_U f, P_U f \rangle = 0 . \end{equation*}

Thus \(f - P_U f\) and \(P_U f\) are orthogonal, and the Pythagorean Theorem (8.9) applied to the decomposition \(f = (f - P_U f) + P_U f\) yields

\begin{equation*} \|f\|^2 = \|f - P_U f\|^2 + \|P_U f\|^2 . \end{equation*}

In particular \(\|P_U f\|^2 \leq \|f\|^2\), and hence

\begin{equation*} \|P_U f\| \leq \|f\| , \end{equation*}

as desired.

Now consider the equality case. The displayed identity shows that \(\|P_U f\| = \|f\|\) if and only if \(\|f - P_U f\|^2 = 0\), which (as \(\|\cdot\|\) is a norm, by 8.19) happens if and only if \(f = P_U f\).

If \(f = P_U f\), then \(f \in U\) because \(P_U f \in U\) by the definition of \(P_U\).

Conversely, suppose \(f \in U\). Then \(\|f - f\| = 0 = \mathrm{distance}(f, U)\), so \(f\) itself is an element of \(U\) closest to \(f\); by the uniqueness in 8.28 the closest element is unique, so \(P_U f = f\).

Hence \(\|P_U f\| = \|f\|\) if and only if \(f \in U\), completing the proof.

Problem (8B.16)

Suppose \(V\) is a Hilbert space and \(P : V \to V\) is a linear map such that \(P^2 = P\) and \(\|Pf\| \leq \|f\|\) for every \(f \in V\). Prove that there exists a closed subspace \(U\) of \(V\) such that \(P = P_U\).

Solution

Let

\begin{equation*} U = \operatorname{range} P . \end{equation*}

Step 1: \(U\) is a closed subspace of \(V\).

Because \(P\) is linear, \(\operatorname{range} P\) is a subspace of \(V\). We claim that

\begin{equation*} \operatorname{range} P = \operatorname{null}(I - P) = \{ f \in V : Pf = f \} . \end{equation*}

Indeed, if \(f \in \operatorname{range} P\), say \(f = Pg\) for some \(g \in V\), then

\begin{equation*} Pf = P(Pg) = P^2 g = Pg = f , \end{equation*}

so \(f \in \operatorname{null}(I - P)\). Conversely, if \(Pf = f\) then \(f \in \operatorname{range} P\). This proves the claim.

The hypothesis \(\|Pf\| \leq \|f\|\) says that \(P\) is a bounded linear map (with \(\|P\| \leq 1\)), hence \(P\) is continuous, hence so is the linear map \(I - P\). Therefore \(\operatorname{null}(I - P) = (I - P)^{-1}(\{0\})\) is the inverse image of the closed set \(\{0\}\) under a continuous function, so it is closed. Thus \(U\) is a closed subspace of \(V\), and consequently \(P_U\) is defined.

Step 2: every element of \(\operatorname{null} P\) is orthogonal to every element of \(U\).

Suppose \(g \in U\) and \(h \in \operatorname{null} P\). For every \(\alpha \in \mathbf{F}\) we have, by Step 1 and by linearity,

\begin{equation*} P(g + \alpha h) = Pg + \alpha Ph = g + 0 = g . \end{equation*}

Hence the hypothesis \(\|P e\| \leq \|e\|\) applied to \(e = g + \alpha h\) gives

\begin{equation*} \|g\| \leq \|g + \alpha h\| \qquad \text{for all } \alpha \in \mathbf{F}. \end{equation*}

Squaring and expanding the inner product (which is linear in the first slot and conjugate linear in the second) gives

\begin{equation*} \|g\|^2 \leq \|g\|^2 + 2 \operatorname{Re}\big( \overline{\alpha} \langle g, h \rangle \big) + |\alpha|^2 \|h\|^2 \end{equation*}

for all \(\alpha \in \mathbf{F}\). Now fix \(t > 0\) and take \(\alpha = -t \langle g, h \rangle\). Then

\begin{equation*} \overline{\alpha} \langle g, h \rangle = -t \, \overline{\langle g, h \rangle} \langle g, h \rangle = -t |\langle g, h \rangle|^2 , \end{equation*}

which is real, and \(|\alpha|^2 = t^2 |\langle g, h \rangle|^2\). The inequality above becomes

\begin{equation*} 0 \leq -2t|\langle g, h \rangle|^2 + t^2 |\langle g, h \rangle|^2 \|h\|^2 , \end{equation*}

that is (after dividing by \(t > 0\)),

\begin{equation*} 2 |\langle g, h \rangle|^2 \leq t \, |\langle g, h \rangle|^2 \|h\|^2 \end{equation*}

for all \(t > 0\). Letting \(t \to 0\) forces \(|\langle g, h \rangle|^2 \leq 0\), and hence

\begin{equation*} \langle g, h \rangle = 0 . \end{equation*}

(This is exactly the implication supplied by Exercise 14 in this section.)

Step 3: \(P = P_U\).

Suppose \(f \in V\). Then \(Pf \in U\). Moreover

\begin{equation*} P(f - Pf) = Pf - P^2 f = Pf - Pf = 0 , \end{equation*}

so \(f - Pf \in \operatorname{null} P\). By Step 2, \(\langle g, f - Pf \rangle = 0\) for every \(g \in U\), and taking complex conjugates gives \(\langle f - Pf, g \rangle = 0\) for every \(g \in U\). Thus \(Pf\) is an element of \(U\) with the property that \(f - Pf\) is orthogonal to every element of \(U\), so 8.37(b) yields

\begin{equation*} Pf = P_U f . \end{equation*}

Since \(f \in V\) was arbitrary, \(P = P_U\) with \(U = \operatorname{range} P\) a closed subspace of \(V\).

Problem (8B.17)

Suppose \(U\) is a subspace of a Hilbert space \(V\). Suppose also that \(W\) is a Banach space and \(S : U \to W\) is a bounded linear map. Prove that there exists a bounded linear map \(T : V \to W\) such that \(T|_U = S\) and \(\|T\| = \|S\|\).

[If \(W = \mathbf{F}\), then this result is just the Hahn-Banach Theorem (6.69) for Hilbert spaces. The result here is stronger because it allows \(W\) to be an arbitrary Banach space instead of requiring \(W\) to be \(\mathbf{F}\). Also, the proof in this Hilbert space context does not require use of Zorn’s Lemma or the Axiom of Choice.]

Solution

If \(U = \{0\}\) then \(S = 0\) and \(\|S\| = 0\), and \(T = 0\) works. So assume \(U \neq \{0\}\).

Step 1: extend \(S\) to the closure of \(U\).

Let \(\overline{U}\) denote the closure of \(U\) in \(V\). By Exercise 12 in Section 6C, \(\overline{U}\) is a subspace of \(V\); it is closed, so \(\overline{U}\) is a closed subspace of the Hilbert space \(V\).

We claim there is a bounded linear map \(\widetilde{S} : \overline{U} \to W\) with \(\widetilde{S}|_U = S\) and \(\|\widetilde{S}\| = \|S\|\); this is Exercise 14 in Section 6C, whose proof we recall. Suppose \(f \in \overline{U}\) and let \(f_1, f_2, \ldots\) be a sequence in \(U\) with \(\lim_{k \to \infty} f_k = f\). Then

\begin{equation*} \|S f_k - S f_j\| = \|S(f_k - f_j)\| \leq \|S\| \, \|f_k - f_j\| , \end{equation*}

and \(f_1, f_2, \ldots\) is Cauchy (being convergent), so \(S f_1, S f_2, \ldots\) is a Cauchy sequence in \(W\). Because \(W\) is a Banach space, this sequence converges; define

\begin{equation*} \widetilde{S} f = \lim_{k \to \infty} S f_k . \end{equation*}

This is well defined: if \(g_1, g_2, \ldots\) is another sequence in \(U\) converging to \(f\), then \(\|S f_k - S g_k\| \leq \|S\| \|f_k - g_k\| \to 0\), so the two limits agree. Taking constant sequences shows \(\widetilde{S}|_U = S\). If \(f, g \in \overline{U}\) and \(\alpha \in \mathbf{F}\), choosing sequences \(f_k \to f\) and \(g_k \to g\) in \(U\) and using \(f_k + \alpha g_k \to f + \alpha g\) together with the linearity of \(S\) and the continuity of the vector space operations shows that \(\widetilde{S}(f + \alpha g) = \widetilde{S} f + \alpha \widetilde{S} g\); thus \(\widetilde{S}\) is linear. Finally, since the norm is a continuous function (Exercise 1 in Section 6C),

\begin{equation*} \|\widetilde{S} f\| = \lim_{k \to \infty} \|S f_k\| \leq \|S\| \lim_{k \to \infty} \|f_k\| = \|S\| \, \|f\| , \end{equation*}

so \(\widetilde{S}\) is bounded with \(\|\widetilde{S}\| \leq \|S\|\). Because \(\widetilde{S}\) extends \(S\), we also have \(\|\widetilde{S}\| \geq \|S\|\). Hence \(\|\widetilde{S}\| = \|S\|\).

Step 2: compose with the orthogonal projection onto \(\overline{U}\).

Because \(\overline{U}\) is a closed subspace of the Hilbert space \(V\), the orthogonal projection \(P_{\overline{U}} : V \to V\) is defined, and it is a linear map by 8.37(c) with \(\operatorname{range} P_{\overline{U}} = \overline{U}\) (8.45(a)). Define

\begin{equation*} T = \widetilde{S} \circ P_{\overline{U}} : V \to W . \end{equation*}

As a composition of linear maps, \(T\) is linear.

If \(f \in U\), then \(f \in \overline{U}\), so \(P_{\overline{U}} f = f\) (the closest element of \(\overline{U}\) to \(f\) is \(f\) itself), and hence

\begin{equation*} Tf = \widetilde{S} f = S f . \end{equation*}

Thus \(T|_U = S\).

For every \(f \in V\) we have, using 8.37(d) (proved in Exercise 15),

\begin{equation*} \|Tf\| = \|\widetilde{S}(P_{\overline{U}} f)\| \leq \|\widetilde{S}\| \, \|P_{\overline{U}} f\| \leq \|\widetilde{S}\| \, \|f\| = \|S\| \, \|f\| . \end{equation*}

Hence \(T\) is bounded with \(\|T\| \leq \|S\|\). On the other hand, \(T\) extends \(S\), so

\begin{equation*} \|T\| = \sup\{\|Tf\| : f \in V, \ \|f\| \leq 1\} \geq \sup\{\|S f\| : f \in U, \ \|f\| \leq 1\} = \|S\| . \end{equation*}

Therefore \(\|T\| = \|S\|\), completing the proof.

Note that the extension was produced by an explicit construction (completion by limits, followed by composition with an orthogonal projection), with no appeal to Zorn’s Lemma.

Problem (8B.18)

Suppose \(U\) and \(W\) are subspaces of a Hilbert space \(V\). Prove that \(\overline{U} = \overline{W}\) if and only if \(U^{\perp} = W^{\perp}\).

Solution

First suppose \(\overline{U} = \overline{W}\). By 8.40(d) applied to \(U\) and then to \(W\),

\begin{equation*} U^{\perp} = \big(\overline{U}\big)^{\perp} = \big(\overline{W}\big)^{\perp} = W^{\perp} . \end{equation*}

Conversely, suppose \(U^{\perp} = W^{\perp}\). Since \(U\) and \(W\) are subspaces of the Hilbert space \(V\), 8.41 applies to each of them and gives

\begin{equation*} \overline{U} = \big(U^{\perp}\big)^{\perp} = \big(W^{\perp}\big)^{\perp} = \overline{W} . \end{equation*}

This completes the proof in both directions.

Two remarks are worth recording. First, the hypothesis that \(V\) is a Hilbert space is needed only for the second implication; the first implication holds in any inner product space, because 8.40(d) does. Second, the statement genuinely requires closures rather than the subspaces themselves: if \(U\) is a dense proper subspace of \(V\) (for example, the subspace of \(\ell^2\) consisting of sequences with only finitely many nonzero coordinates) and \(W = V\), then \(U^{\perp} = \{0\} = W^{\perp}\) by 8.42, yet \(U \neq W\).

In particular, taking \(U\) and \(W\) to be closed subspaces, we get: two closed subspaces of a Hilbert space are equal if and only if their orthogonal complements are equal.

Problem (8B.19)

Suppose \(U\) and \(W\) are closed subspaces of a Hilbert space. Prove that \(P_U P_W = 0\) if and only if \(\langle f, g \rangle = 0\) for all \(f \in U\) and all \(g \in W\).

Solution

Call the Hilbert space \(V\). Because \(U\) and \(W\) are closed subspaces of \(V\), the orthogonal projections \(P_U\) and \(P_W\) are defined and are linear maps from \(V\) to \(V\) (8.37(c)). Throughout we use 8.45(a), which states that

\begin{equation*} \operatorname{null} P_U = U^{\perp} \quad \text{and} \quad \operatorname{range} P_W = W . \end{equation*}

First suppose that \(\langle f, g \rangle = 0\) for all \(f \in U\) and all \(g \in W\). By the definition of orthogonal complement (8.38), this says precisely that

\begin{equation*} W \subseteq U^{\perp} . \end{equation*}

Now let \(h \in V\). Then \(P_W h \in W \subseteq U^{\perp} = \operatorname{null} P_U\), and hence

\begin{equation*} (P_U P_W)(h) = P_U (P_W h) = 0 . \end{equation*}

Since \(h \in V\) was arbitrary, \(P_U P_W = 0\).

Conversely, suppose \(P_U P_W = 0\). Let \(g \in W\). Because \(g \in W\) and \(W\) is a closed subspace, \(g\) is the element of \(W\) closest to \(g\), so \(P_W g = g\). Therefore

\begin{equation*} P_U g = P_U (P_W g) = (P_U P_W)(g) = 0 , \end{equation*}

so \(g \in \operatorname{null} P_U = U^{\perp}\). By the definition of \(U^{\perp}\), this means \(\langle f, g \rangle = 0\) for every \(f \in U\). As \(g \in W\) was arbitrary, we conclude that \(\langle f, g \rangle = 0\) for all \(f \in U\) and all \(g \in W\), as desired.

Note that the condition is symmetric in \(U\) and \(W\) (because \(\langle g, f \rangle = \overline{\langle f, g \rangle}\)), so the result also shows that \(P_U P_W = 0\) if and only if \(P_W P_U = 0\).

Problem (8B.20)

Verify the assertions in Example 8.46. That example states: suppose \(U\) is the closed subspace of \(L^2(\mathbb{R})\) defined by

\begin{equation*} U = \{ f \in L^2(\mathbb{R}) : f(x) = 0 \text{ for almost every } x < 0 \} . \end{equation*}

Then

\begin{equation*} U^{\perp} = \{ g \in L^2(\mathbb{R}) : g(x) = 0 \text{ for almost every } x \geq 0 \} , \end{equation*}

and if \(h \in L^2(\mathbb{R})\), then

\begin{equation*} P_U h = h \chi_{[0, \infty)} \quad \text{and} \quad P_{U^{\perp}} h = h \chi_{(-\infty, 0)} . \end{equation*}

Thus \(P_{U^{\perp}} h = h(1 - \chi_{[0, \infty)}) = (I - P_U) h\) and hence \(P_{U^{\perp}} = I - P_U\), as asserted in 8.45(c).

Solution

Throughout, \(L^2(\mathbb{R})\) means \(L^2\) of Lebesgue measure on \(\mathbb{R}\), with inner product \(\langle f, g \rangle = \int f \overline{g} \, d\lambda\), and its elements are equivalence classes of functions that agree almost everywhere; the conditions defining \(U\) and the set below are therefore well posed. Write

\begin{equation*} W = \{ g \in L^2(\mathbb{R}) : g(x) = 0 \text{ for almost every } x \geq 0 \} . \end{equation*}

Step 1: \(U\) is a closed subspace of \(L^2(\mathbb{R})\).

If \(f_1, f_2 \in U\) and \(\alpha \in \mathbf{F}\), then \(f_1 + \alpha f_2\) vanishes almost everywhere on \((-\infty, 0)\) (the union of two sets of measure zero has measure zero), so \(U\) is a subspace. To see that \(U\) is closed, suppose \(f_1, f_2, \ldots\) is a sequence in \(U\) converging in \(L^2(\mathbb{R})\) to some \(f \in L^2(\mathbb{R})\). For each \(k\) we have \(f_k = 0\) almost everywhere on \((-\infty, 0)\), and hence

\begin{equation*} \int_{(-\infty, 0)} |f|^2 \, d\lambda = \int_{(-\infty, 0)} |f - f_k|^2 \, d\lambda \leq \|f - f_k\|^2 \longrightarrow 0 . \end{equation*}

Thus \(\int_{(-\infty, 0)} |f|^2 \, d\lambda = 0\), which implies \(f = 0\) almost everywhere on \((-\infty, 0)\), so \(f \in U\). Hence \(U\) is closed, and therefore \(P_U\) is defined. Similarly \(W\) is a closed subspace (same argument with \([0, \infty)\) in place of \((-\infty, 0)\)).

Step 2: \(U^{\perp} = W\).

Suppose \(g \in W\) and \(f \in U\). Then \(f \overline{g} = 0\) almost everywhere: on \((-\infty, 0)\) we have \(f = 0\) almost everywhere, and on \([0, \infty)\) we have \(g = 0\) almost everywhere. Hence

\begin{equation*} \langle f, g \rangle = \int_{\mathbb{R}} f \overline{g} \, d\lambda = 0 . \end{equation*}

As this holds for all \(f \in U\), we get \(g \in U^{\perp}\); thus \(W \subseteq U^{\perp}\).

Conversely, suppose \(g \in U^{\perp}\). Put \(f = g \chi_{[0, \infty)}\). Then \(|f| \leq |g|\) pointwise, so \(f \in L^2(\mathbb{R})\), and \(f = 0\) on \((-\infty, 0)\), so \(f \in U\). Hence

\begin{equation*} 0 = \langle f, g \rangle = \int_{\mathbb{R}} g \chi_{[0, \infty)} \overline{g} \, d\lambda = \int_{[0, \infty)} |g|^2 \, d\lambda . \end{equation*}

Therefore \(g = 0\) almost everywhere on \([0, \infty)\), so \(g \in W\). Thus \(U^{\perp} \subseteq W\), and combining the two inclusions gives \(U^{\perp} = W\).

Step 3: \(P_U h = h \chi_{[0, \infty)}\) for every \(h \in L^2(\mathbb{R})\).

Let \(h \in L^2(\mathbb{R})\) and set \(h_1 = h \chi_{[0, \infty)}\). Since \(|h_1| \leq |h|\) we have \(h_1 \in L^2(\mathbb{R})\), and \(h_1 = 0\) on \((-\infty, 0)\), so \(h_1 \in U\). Moreover

\begin{equation*} h - h_1 = h\big(1 - \chi_{[0, \infty)}\big) = h \chi_{(-\infty, 0)} , \end{equation*}

which lies in \(L^2(\mathbb{R})\) and vanishes on \([0, \infty)\), so \(h - h_1 \in W = U^{\perp}\) by Step 2. In particular \(h - h_1\) is orthogonal to every element of \(U\). Since \(h_1 \in U\), part (b) of 8.37 gives

\begin{equation*} P_U h = h_1 = h \chi_{[0, \infty)} . \end{equation*}

Step 4: \(P_{U^{\perp}} h = h \chi_{(-\infty, 0)}\) for every \(h \in L^2(\mathbb{R})\).

By 8.40(a), \(U^{\perp}\) is a closed subspace, so \(P_{U^{\perp}}\) is defined. Set \(h_2 = h \chi_{(-\infty, 0)}\); by Step 2 (and the argument in Step 3) we have \(h_2 \in W = U^{\perp}\). Also

\begin{equation*} h - h_2 = h \chi_{[0, \infty)} \in U , \end{equation*}

and every element of \(U\) is orthogonal to every element of \(U^{\perp}\) by the definition of \(U^{\perp}\) (equivalently, \(U \subseteq (U^{\perp})^{\perp}\) by 8.40(e)). Hence \(h - h_2\) is orthogonal to every element of \(U^{\perp}\). Applying 8.37(b) with \(U^{\perp}\) in place of \(U\) gives

\begin{equation*} P_{U^{\perp}} h = h_2 = h \chi_{(-\infty, 0)} . \end{equation*}

Step 5: \(P_{U^{\perp}} = I - P_U\).

For \(h \in L^2(\mathbb{R})\), Steps 3 and 4 give

\begin{equation*} P_{U^{\perp}} h = h \chi_{(-\infty, 0)} = h\big(1 - \chi_{[0, \infty)}\big) = h - h\chi_{[0, \infty)} = h - P_U h = (I - P_U) h . \end{equation*}

Since \(h\) was arbitrary, \(P_{U^{\perp}} = I - P_U\), confirming 8.45(c) in this example. All assertions of Example 8.46 are verified.

Problem (8B.21)

Show that every inner product space is a subspace of some Hilbert space.

Hint: See Exercise 13 in Section 6C.

Solution

Suppose \(V\) is an inner product space, with inner product \(\langle \cdot, \cdot \rangle_V\) and with the norm \(\|f\| = \sqrt{\langle f, f \rangle_V}\) (a norm by 8.19). We must produce a Hilbert space \(B\) such that \(V\) is a subspace of \(B\) and the inner product of \(B\) restricted to \(V\) equals \(\langle \cdot, \cdot \rangle_V\).

Step 1: a Banach space containing \(V\) as a dense subspace.

By Exercise 13 in Section 6C (whose construction starts from the completion of the metric space \(V\), as in Exercise 16 in Section 6A), the normed vector space \(V\) is a subspace of some Banach space \(C\) whose norm restricted to \(V\) is the norm of \(V\). We denote the norm of \(C\) again by \(\|\cdot\|\).

We may arrange that \(V\) is dense, without relying on any feature of that construction beyond its statement. Let

\begin{equation*} B = \overline{V} , \end{equation*}

the closure of \(V\) in \(C\). By Exercise 12 in Section 6C, \(B\) is a subspace of \(C\), and it is closed. A closed subset of a complete metric space is complete: if \(b_1, b_2, \ldots\) is a Cauchy sequence in \(B\), then it converges to some \(b \in C\) because \(C\) is a Banach space, and \(b \in \overline{B} = B\) because \(B\) is closed. Hence \(B\) is a Banach space, its norm restricted to \(V\) is the norm of \(V\), and by construction \(V\) is dense in \(B\): every element of \(B\) is the limit of a sequence in \(V\).

Step 2: extend the inner product to \(B\) by continuity.

Let \(f, g \in B\) and choose sequences \(f_1, f_2, \ldots\) and \(g_1, g_2, \ldots\) in \(V\) with \(f_k \to f\) and \(g_k \to g\) in \(B\). Convergent sequences are bounded (this follows from Exercise 4 in Section 6C, since convergent sequences are Cauchy), so there is \(M > 0\) with \(\|f_k\| \leq M\) and \(\|g_k\| \leq M\) for all \(k\). For all \(j, k\), the Cauchy-Schwarz inequality (8.11) in \(V\) gives

\begin{equation*} \begin{aligned} \big| \langle f_k, g_k \rangle_V - \langle f_j, g_j \rangle_V \big| &= \big| \langle f_k - f_j, g_k \rangle_V + \langle f_j, g_k - g_j \rangle_V \big| \\ &\leq \|f_k - f_j\| \, \|g_k\| + \|f_j\| \, \|g_k - g_j\| \\ &\leq M \big( \|f_k - f_j\| + \|g_k - g_j\| \big) . \end{aligned} \end{equation*}

Because \((f_k)\) and \((g_k)\) are Cauchy, the right side tends to \(0\) as \(j, k \to \infty\). Thus \(\langle f_1, g_1 \rangle_V, \langle f_2, g_2 \rangle_V, \ldots\) is a Cauchy sequence in \(\mathbf{F}\), hence converges.

The limit does not depend on the chosen sequences: if also \(f_k^{\prime} \to f\) and \(g_k^{\prime} \to g\) with \(f_k^{\prime}, g_k^{\prime} \in V\), then the same estimate (with a common bound \(M\)) gives

\begin{equation*} \big| \langle f_k, g_k \rangle_V - \langle f_k^{\prime}, g_k^{\prime} \rangle_V \big| \leq M \big( \|f_k - f_k^{\prime}\| + \|g_k - g_k^{\prime}\| \big) \longrightarrow 0 . \end{equation*}

So we may define

\begin{equation*} \langle f, g \rangle = \lim_{k \to \infty} \langle f_k, g_k \rangle_V \qquad (f, g \in B). \end{equation*}

Taking constant sequences shows that \(\langle f, g \rangle = \langle f, g \rangle_V\) whenever \(f, g \in V\).

Step 3: the extension is an inner product on \(B\) whose norm is the norm of \(B\).

Suppose \(f, f^{\prime}, g \in B\) and \(\alpha \in \mathbf{F}\); choose sequences \(f_k \to f\), \(f_k^{\prime} \to f^{\prime}\), \(g_k \to g\) in \(V\). Since addition and scalar multiplication are continuous on a normed vector space, \(f_k + \alpha f_k^{\prime} \to f + \alpha f^{\prime}\), and these vectors lie in \(V\). Hence

\begin{equation*} \langle f + \alpha f^{\prime}, g \rangle = \lim_{k \to \infty} \langle f_k + \alpha f_k^{\prime}, g_k \rangle_V = \lim_{k \to \infty} \big( \langle f_k, g_k \rangle_V + \alpha \langle f_k^{\prime}, g_k \rangle_V \big) = \langle f, g \rangle + \alpha \langle f^{\prime}, g \rangle , \end{equation*}

so \(\langle \cdot, \cdot \rangle\) is linear in its first slot. Likewise, conjugation is continuous on \(\mathbf{F}\), so

\begin{equation*} \langle g, f \rangle = \lim_{k \to \infty} \langle g_k, f_k \rangle_V = \lim_{k \to \infty} \overline{\langle f_k, g_k \rangle_V} = \overline{\langle f, g \rangle} , \end{equation*}

giving conjugate symmetry.

For positivity, take \(g = f\) and \(g_k = f_k\). Because the norm is a continuous function (Exercise 1 in Section 6C), \(\|f_k\| \to \|f\|\), and therefore

\begin{equation*} \langle f, f \rangle = \lim_{k \to \infty} \langle f_k, f_k \rangle_V = \lim_{k \to \infty} \|f_k\|^2 = \|f\|^2 . \end{equation*}

Thus \(\langle f, f \rangle \geq 0\), with \(\langle f, f \rangle = 0\) if and only if \(\|f\| = 0\), that is, if and only if \(f = 0\). Hence \(\langle \cdot, \cdot \rangle\) is an inner product on \(B\), and the displayed identity says that the norm it induces is exactly the norm of \(B\):

\begin{equation*} \sqrt{\langle f, f \rangle} = \|f\| \qquad \text{for all } f \in B . \end{equation*}

Step 4: conclusion.

\(B\) is an inner product space whose associated norm is its original Banach space norm, so \(B\) is complete with respect to the norm coming from its inner product; that is, \(B\) is a Hilbert space (8.21). By Step 1, \(V\) is a subspace of \(B\), and by the last sentence of Step 2, the inner product of \(B\) restricted to \(V\) is the original inner product of \(V\). Therefore every inner product space is a subspace of some Hilbert space.

A final remark: the inner product constructed in Step 2 is the only possible extension. Indeed, if \(\langle \cdot, \cdot \rangle^{\prime}\) is any inner product on \(B\) that induces the norm of \(B\) and agrees with \(\langle \cdot, \cdot \rangle_V\) on \(V\), then for \(f_k \to f\) and \(g_k \to g\) with \(f_k, g_k \in V\) the Cauchy-Schwarz inequality (8.11) applied in \(B\) gives

\begin{equation*} \big| \langle f, g \rangle^{\prime} - \langle f_k, g_k \rangle^{\prime} \big| \leq \|f - f_k\| \, \|g\| + \|f_k\| \, \|g - g_k\| \longrightarrow 0 , \end{equation*}

so \(\langle f, g \rangle^{\prime} = \lim_{k \to \infty} \langle f_k, g_k \rangle_V = \langle f, g \rangle\).

Problem (8B.22)

Prove that if \(V\) is a Hilbert space and \(T : V \to V\) is a bounded linear map such that the dimension of \(\operatorname{range} T\) is 1, then there exist \(g, h \in V\) such that

\begin{equation*} T f = \langle f , g \rangle h \end{equation*}

for all \(f \in V\).

Solution

Because \(\operatorname{range} T\) is a one-dimensional subspace of \(V\), there exists \(h \in V\) with \(h \neq 0\) such that

\begin{equation*} \operatorname{range} T = \{ \alpha h : \alpha \in \mathbf{F} \}. \end{equation*}

Hence for each \(f \in V\) there exists \(\alpha \in \mathbf{F}\) with \(T f = \alpha h\), and this scalar \(\alpha\) is unique: if \(\alpha h = \beta h\) then \((\alpha - \beta) h = 0\), which forces \(\alpha = \beta\) because \(h \neq 0\). Thus we may define a function \(\varphi : V \to \mathbf{F}\) by declaring \(\varphi(f)\) to be the unique scalar such that

\begin{equation*} T f = \varphi(f) h . \end{equation*}

\(\varphi\) is linear. Suppose \(f_1, f_2 \in V\) and \(\alpha \in \mathbf{F}\). Then

\begin{equation*} \varphi(f_1 + f_2) h = T(f_1 + f_2) = T f_1 + T f_2 = \bigl( \varphi(f_1) + \varphi(f_2) \bigr) h \end{equation*}

and

\begin{equation*} \varphi(\alpha f_1) h = T(\alpha f_1) = \alpha T f_1 = \bigl( \alpha \varphi(f_1) \bigr) h . \end{equation*}

The uniqueness noted above now gives \(\varphi(f_1 + f_2) = \varphi(f_1) + \varphi(f_2)\) and \(\varphi(\alpha f_1) = \alpha \varphi(f_1)\). Thus \(\varphi\) is a linear functional on \(V\).

\(\varphi\) is bounded. For each \(f \in V\) we have

\begin{equation*} |\varphi(f)| \, \|h\| = \|\varphi(f) h\| = \|T f\| \le \|T\| \, \|f\| , \end{equation*}

and therefore

\begin{equation*} |\varphi(f)| \le \frac{\|T\|}{\|h\|} \, \|f\| , \end{equation*}

which shows that \(\varphi\) is a bounded linear functional on \(V\) (with \(\|\varphi\| \le \|T\| / \|h\|\)).

Apply the Riesz Representation Theorem. Because \(V\) is a Hilbert space and \(\varphi\) is a bounded linear functional on \(V\), the Riesz Representation Theorem (8.47) provides \(g \in V\) such that

\begin{equation*} \varphi(f) = \langle f , g \rangle \quad \text{for all } f \in V . \end{equation*}

Combining the last equation with the definition of \(\varphi\), we conclude that

\begin{equation*} T f = \varphi(f) h = \langle f , g \rangle h \end{equation*}

for all \(f \in V\), as desired.

Note that the hypothesis that \(V\) is complete enters only through the appeal to 8.47; every step before that is valid in an arbitrary inner product space.

Problem (8B.23)

(a) Give an example of a Banach space \(V\) and a bounded linear functional \(\varphi\) on \(V\) such that \(|\varphi(f)| < \|\varphi\| \, \|f\|\) for all \(f \in V \setminus \{0\}\).

(b) Show there does not exist an example in part (a) where \(V\) is a Hilbert space.

Solution

(a) Let \(V = \ell^1\), which is \(L^1(\mu)\) for \(\mu\) counting measure on \(\mathbf{Z}^+\) and hence is a Banach space by 7.24. Define \(\varphi : \ell^1 \to \mathbf{F}\) by

\begin{equation*} \varphi(a) = \sum_{k=1}^{\infty} \Bigl( 1 - \frac{1}{k} \Bigr) a_k \qquad \text{for } a = (a_1, a_2, \dots) \in \ell^1 . \end{equation*}

The series converges absolutely because \(|(1 - 1/k) a_k| \le |a_k|\) and \(\sum_k |a_k| < \infty\); linearity of \(\varphi\) is clear.

The norm of \(\varphi\) is 1. For every \(a \in \ell^1\),

\begin{equation*} |\varphi(a)| \le \sum_{k=1}^{\infty} \Bigl( 1 - \frac{1}{k} \Bigr) |a_k| \le \sum_{k=1}^{\infty} |a_k| = \|a\|_1 , \end{equation*}

so \(\|\varphi\| \le 1\). In the other direction, let \(e_k \in \ell^1\) be the sequence whose \(k^{\text{th}}\) coordinate is 1 and whose other coordinates are 0. Then \(\|e_k\|_1 = 1\) and \(\varphi(e_k) = 1 - 1/k\), so \(\|\varphi\| \ge 1 - 1/k\) for every \(k \in \mathbf{Z}^+\); letting \(k \to \infty\) gives \(\|\varphi\| \ge 1\). Hence \(\|\varphi\| = 1\).

The inequality is strict. Suppose \(a \in \ell^1\) with \(a \neq 0\). Choose \(m \in \mathbf{Z}^+\) with \(a_m \neq 0\). Then

\begin{equation*} |\varphi(a)| \le \sum_{k=1}^{\infty} \Bigl( 1 - \frac{1}{k} \Bigr) |a_k| = \|a\|_1 - \sum_{k=1}^{\infty} \frac{|a_k|}{k} \le \|a\|_1 - \frac{|a_m|}{m} < \|a\|_1 = \|\varphi\| \, \|a\|_1 , \end{equation*}

where every quantity in sight is finite because \(a \in \ell^1\). Thus \(|\varphi(a)| < \|\varphi\| \, \|a\|\) for all \(a \in \ell^1 \setminus \{0\}\), as required.

(b) Suppose \(V\) is a Hilbert space with \(V \neq \{0\}\) and \(\varphi\) is a bounded linear functional on \(V\). We show that the strict inequality of (a) must fail for some \(f \in V \setminus \{0\}\). (If \(V = \{0\}\) then \(V \setminus \{0\} = \emptyset\) and the condition holds vacuously, so the nontrivial case is the one just described.)

By the Riesz Representation Theorem (8.47), there exists \(h \in V\) such that

\begin{equation*} \varphi(f) = \langle f , h \rangle \quad \text{for all } f \in V , \qquad \text{and} \qquad \|\varphi\| = \|h\| . \end{equation*}

If \(h \neq 0\), take \(f = h\). Then \(f \in V \setminus \{0\}\) and

\begin{equation*} |\varphi(f)| = |\langle h , h \rangle| = \|h\|^2 = \|\varphi\| \, \|h\| = \|\varphi\| \, \|f\| , \end{equation*}

so the strict inequality fails at \(f = h\).

If \(h = 0\), then \(\varphi = 0\) and \(\|\varphi\| = 0\). Because \(V \neq \{0\}\) we may choose any \(f \in V \setminus \{0\}\), and then

\begin{equation*} |\varphi(f)| = 0 = \|\varphi\| \, \|f\| , \end{equation*}

so again the strict inequality fails.

In either case the strict inequality fails at some \(f \in V \setminus \{0\}\), so no Hilbert space \(V \neq \{0\}\) furnishes an example as in (a). (The excluded space \(V = \{0\}\) is degenerate: there \(V \setminus \{0\}\) is empty and the condition in (a) is vacuously satisfied by the only linear functional on \(V\).) The content of the argument is that on a Hilbert space every bounded linear functional attains its norm, namely at the vector \(h\) representing it.

Problem (8B.24)

(a) Suppose \(\varphi\) and \(\psi\) are bounded linear functionals on a Hilbert space \(V\) such that \(\|\varphi + \psi\| = \|\varphi\| + \|\psi\|\). Prove that one of \(\varphi, \psi\) is a scalar multiple of the other.

(b) Give an example to show that (a) can fail if the hypothesis that \(V\) is a Hilbert space is replaced by the hypothesis that \(V\) is a Banach space.

Solution

(a) By the Riesz Representation Theorem (8.47) there exist \(g, h \in V\) such that

\begin{equation*} \varphi(f) = \langle f , g \rangle \quad \text{and} \quad \psi(f) = \langle f , h \rangle \qquad \text{for all } f \in V , \end{equation*}

and moreover \(\|\varphi\| = \|g\|\) and \(\|\psi\| = \|h\|\). Because

\begin{equation*} (\varphi + \psi)(f) = \langle f , g \rangle + \langle f , h \rangle = \langle f , g + h \rangle \qquad \text{for all } f \in V , \end{equation*}

the uniqueness assertion in 8.47 shows that \(g + h\) is the vector representing \(\varphi + \psi\); hence \(\|\varphi + \psi\| = \|g + h\|\). The hypothesis therefore becomes

\begin{equation*} \|g + h\| = \|g\| + \|h\| . \end{equation*}

Squaring both sides and expanding the left side using the properties of the inner product gives

\begin{equation*} \|g\|^2 + 2 \operatorname{Re} \langle g , h \rangle + \|h\|^2 = \|g\|^2 + 2 \|g\| \, \|h\| + \|h\|^2 , \end{equation*}

and hence

\begin{equation*} \operatorname{Re} \langle g , h \rangle = \|g\| \, \|h\| . \end{equation*}

Because \(\operatorname{Re} \langle g, h \rangle \le |\langle g , h \rangle| \le \|g\| \, \|h\|\) by the Cauchy–Schwarz inequality (8.11), both of these inequalities are equalities. In particular

\begin{equation*} |\langle g , h \rangle| = \|g\| \, \|h\| , \end{equation*}

which is the equality case of 8.11; hence one of \(g, h\) is a scalar multiple of the other.

For completeness here is the direct verification of that last step. If \(h = 0\), then \(\psi = 0 = 0 \cdot \varphi\) and we are done. So assume \(h \neq 0\) and set

\begin{equation*} u = g - \frac{\langle g , h \rangle}{\|h\|^2} \, h . \end{equation*}

Then \(\langle u , h \rangle = \langle g, h \rangle - \frac{\langle g,h\rangle}{\|h\|^2}\|h\|^2 = 0\), so \(u\) is orthogonal to \(h\), and the Pythagorean Theorem (8.9) gives

\begin{equation*} \|g\|^2 = \|u\|^2 + \frac{|\langle g , h \rangle|^2}{\|h\|^2} = \|u\|^2 + \|g\|^2 , \end{equation*}

using \(|\langle g,h\rangle| = \|g\|\,\|h\|\). Hence \(\|u\| = 0\), so

\begin{equation*} g = \frac{\langle g , h \rangle}{\|h\|^2} \, h = \frac{\|g\|}{\|h\|} \, h , \end{equation*}

where the last equality uses \(\langle g , h \rangle = \operatorname{Re}\langle g,h\rangle = \|g\| \, \|h\|\) (the imaginary part vanishes because the modulus equals the real part).

Writing \(c = \|g\| / \|h\|\), which is a nonnegative real number, we get for all \(f \in V\)

\begin{equation*} \varphi(f) = \langle f , g \rangle = \langle f , c h \rangle = \overline{c} \, \langle f , h \rangle = c \, \psi(f) , \end{equation*}

since \(c\) is real. Thus \(\varphi = c \psi\), so \(\varphi\) is a scalar multiple of \(\psi\), completing the proof of (a).

(b) Let \(V = \mathbf{F}^2\) with the norm

\begin{equation*} \|(a_1, a_2)\| = \max \{ |a_1| , |a_2| \} . \end{equation*}

Convergence with respect to this norm is exactly coordinatewise convergence, and a sequence is Cauchy in this norm exactly when each coordinate sequence is Cauchy in \(\mathbf{F}\); because \(\mathbf{F}\) is complete, \(V\) is a Banach space.

Define bounded linear functionals \(\varphi, \psi\) on \(V\) by

\begin{equation*} \varphi(a_1, a_2) = a_1 , \qquad \psi(a_1, a_2) = a_2 . \end{equation*}

Then \(|\varphi(a_1,a_2)| = |a_1| \le \|(a_1,a_2)\|\) with equality at \((1,0)\), so \(\|\varphi\| = 1\); similarly \(\|\psi\| = 1\). Also \((\varphi + \psi)(a_1, a_2) = a_1 + a_2\), so

\begin{equation*} |(\varphi + \psi)(a_1, a_2)| \le |a_1| + |a_2| \le 2 \|(a_1, a_2)\| , \end{equation*}

giving \(\|\varphi + \psi\| \le 2\), while \((\varphi + \psi)(1,1) = 2\) and \(\|(1,1)\| = 1\) give \(\|\varphi + \psi\| \ge 2\). Hence

\begin{equation*} \|\varphi + \psi\| = 2 = \|\varphi\| + \|\psi\| . \end{equation*}

However neither of \(\varphi, \psi\) is a scalar multiple of the other: if \(\varphi = c \psi\) then evaluating at \((1,0)\) gives \(1 = c \cdot 0 = 0\), a contradiction, and if \(\psi = c \varphi\) then evaluating at \((0,1)\) gives \(1 = 0\), again a contradiction. Thus (a) fails for this Banach space.

Problem (8B.25)

(a) Suppose that \(\mu\) is a finite measure, \(1 \le p \le 2\), and \(\varphi\) is a bounded linear functional on \(L^p(\mu)\). Prove that there exists \(h \in L^{p^{\prime}}(\mu)\) such that \(\varphi(f) = \int f h \, d\mu\) for every \(f \in L^p(\mu)\).

(b) Same as (a), but with the hypothesis that \(\mu\) is a finite measure replaced by the hypothesis that \(\mu\) is a measure, and assume that \(1 < p \le 2\).

[See 7.25, which along with this exercise shows that we can identify the dual of \(L^p(\mu)\) with \(L^{p^{\prime}}(\mu)\) for \(1 < p \le 2\). See 9.42 for an extension to all \(p \in (1, \infty)\).]

Solution

Throughout, \((X, \mathcal{S}, \mu)\) is the underlying measure space, \(p^{\prime}\) denotes the exponent with \(\frac{1}{p} + \frac{1}{p^{\prime}} = 1\), and for a measurable \(h\) we write \(\sigma\) for the measurable function defined by \(\sigma = \overline{h} / |h|\) on \(\{ h \neq 0 \}\) and \(\sigma = 0\) elsewhere; thus \(|\sigma| \le 1\) and \(\sigma h = |h|\).

(a) If \(\mu(X) = 0\) then every \(L^q(\mu)\) equals \(\{0\}\) and \(h = 0\) works, so assume \(\mu(X) > 0\). Note \(1 \le p \le 2\) forces \(2 \le p^{\prime} \le \infty\).

Step 1: \(L^2(\mu) \subseteq L^p(\mu)\), continuously. Let \(f \in L^2(\mu)\). If \(p = 2\) there is nothing to prove. If \(p < 2\), apply Hölder’s inequality (7.9) to the functions \(|f|^p\) and \(1\) with the exponents \(\frac{2}{p}\) and \(\frac{2}{2-p}\) (which are conjugate):

\begin{equation*} \int |f|^p \, d\mu \le \Bigl( \int |f|^2 \, d\mu \Bigr)^{p/2} \mu(X)^{(2-p)/2} . \end{equation*}

Raising to the power \(1/p\) gives

\begin{equation*} \|f\|_p \le c \, \|f\|_2 , \qquad \text{where } c = \mu(X)^{\frac{1}{p} - \frac{1}{2}} \in (0, \infty) . \end{equation*}

In particular \(L^2(\mu) \subseteq L^p(\mu)\).

Step 2: represent \(\varphi\) on \(L^2(\mu)\). By Step 1 the restriction of \(\varphi\) to \(L^2(\mu)\) is a linear functional satisfying

\begin{equation*} |\varphi(f)| \le \|\varphi\| \, \|f\|_p \le c \, \|\varphi\| \, \|f\|_2 \qquad \text{for all } f \in L^2(\mu) , \end{equation*}

so it is a bounded linear functional on the Hilbert space \(L^2(\mu)\) (8.22). By the Riesz Representation Theorem (8.47) there is \(g \in L^2(\mu)\) with \(\varphi(f) = \langle f , g \rangle = \int f \overline{g} \, d\mu\) for all \(f \in L^2(\mu)\). Put \(h = \overline{g} \in L^2(\mu)\), so that

\begin{equation*} \varphi(f) = \int f h \, d\mu \qquad \text{for all } f \in L^2(\mu) . \end{equation*}

(Here \(\langle f , g \rangle = \int f \overline{g} \, d\mu\) is the usual inner product on \(L^2(\mu)\), the one under which \(L^2(\mu)\) is a Hilbert space in 8.22; the integral is finite by the Cauchy–Schwarz inequality 8.11.)

Step 3: \(h \in L^{p^{\prime}}(\mu)\) with \(\|h\|_{p^{\prime}} \le \|\varphi\|\). Note first that because \(\mu\) is finite, every bounded measurable function lies in \(L^2(\mu) \cap L^p(\mu)\).

Suppose first \(1 < p \le 2\), so \(2 \le p^{\prime} < \infty\). Since \(h \in L^2(\mu)\) we have \(|h| < \infty\) almost everywhere; discarding a set of measure 0, assume \(|h| < \infty\) everywhere. Let \(E_n = \{ x \in X : |h(x)| \le n \}\) and set

\begin{equation*} f_n = \chi_{E_n} |h|^{p^{\prime} - 1} \sigma . \end{equation*}

Then \(|f_n| \le n^{p^{\prime}-1}\), so \(f_n \in L^2(\mu)\), and \(f_n h = \chi_{E_n} |h|^{p^{\prime}}\). Writing \(I_n = \int_{E_n} |h|^{p^{\prime}} \, d\mu\) (which is finite, being at most \(n^{p^{\prime}} \mu(X)\)), Step 2 gives \(\varphi(f_n) = I_n\). Because \((p^{\prime} - 1) p = p^{\prime}\), we also have

\begin{equation*} \|f_n\|_p^p = \int \chi_{E_n} |h|^{(p^{\prime}-1)p} \, d\mu = I_n . \end{equation*}

Hence

\begin{equation*} I_n = |\varphi(f_n)| \le \|\varphi\| \, \|f_n\|_p = \|\varphi\| \, I_n^{1/p} . \end{equation*}

If \(I_n > 0\) we may divide by \(I_n^{1/p}\) to get \(I_n^{1 - 1/p} = I_n^{1/p^{\prime}} \le \|\varphi\|\), and this inequality holds trivially when \(I_n = 0\). The functions \(\chi_{E_n} |h|^{p^{\prime}}\) increase pointwise to \(|h|^{p^{\prime}}\), so the Monotone Convergence Theorem (3.11) gives

\begin{equation*} \int |h|^{p^{\prime}} \, d\mu = \lim_{n \to \infty} I_n \le \|\varphi\|^{p^{\prime}} , \end{equation*}

that is, \(h \in L^{p^{\prime}}(\mu)\) with \(\|h\|_{p^{\prime}} \le \|\varphi\|\).

Now suppose \(p = 1\), so \(p^{\prime} = \infty\). Let \(\varepsilon > 0\) and let \(E = \{ x \in X : |h(x)| > \|\varphi\| + \varepsilon \}\). The function \(f = \chi_E \sigma\) is bounded, hence lies in \(L^2(\mu) \cap L^1(\mu)\), and \(f h = \chi_E |h|\), while \(\|f\|_1 = \mu(E)\) because \(|\sigma| = 1\) on \(E\). Therefore

\begin{equation*} (\|\varphi\| + \varepsilon) \mu(E) \le \int_E |h| \, d\mu = \varphi(f) \le \|\varphi\| \, \mu(E) . \end{equation*}

Since \(\mu(E) \le \mu(X) < \infty\), this forces \(\varepsilon \mu(E) \le 0\), so \(\mu(E) = 0\). As \(\varepsilon > 0\) was arbitrary, \(|h| \le \|\varphi\|\) almost everywhere; thus \(h \in L^{\infty}(\mu)\) with \(\|h\|_{\infty} \le \|\varphi\|\).

Step 4: the formula holds on all of \(L^p(\mu)\). Define \(\psi : L^p(\mu) \to \mathbf{F}\) by \(\psi(f) = \int f h \, d\mu\). By Hölder’s inequality (7.9), \(f h \in L^1(\mu)\) and \(|\psi(f)| \le \|h\|_{p^{\prime}} \|f\|_p\), so \(\psi\) is a bounded linear functional on \(L^p(\mu)\). By Step 2, \(\varphi\) and \(\psi\) agree on \(L^2(\mu)\).

We claim \(L^2(\mu)\) is dense in \(L^p(\mu)\). Indeed, let \(f \in L^p(\mu)\) and set \(f_n = f \chi_{\{|f| \le n\}}\). Each \(f_n\) is bounded, hence in \(L^2(\mu)\) because \(\mu\) is finite, and

\begin{equation*} \|f - f_n\|_p^p = \int_{\{|f| > n\}} |f|^p \, d\mu \longrightarrow 0 \end{equation*}

by the Dominated Convergence Theorem (3.31): the integrands \(\chi_{\{|f| > n\}} |f|^p\) are dominated by \(|f|^p \in L^1(\mu)\) and converge pointwise to 0 (because \(|f| < \infty\) almost everywhere). Thus \(f_n \to f\) in \(L^p(\mu)\), proving the claim. Note \(p < \infty\) here, which is why this argument works.

The bounded linear functional \(\varphi - \psi\) is continuous and vanishes on the dense subset \(L^2(\mu)\) of \(L^p(\mu)\), hence vanishes identically. Therefore

\begin{equation*} \varphi(f) = \int f h \, d\mu \qquad \text{for all } f \in L^p(\mu) , \end{equation*}

completing the proof of (a).

(b) Now \(\mu\) is an arbitrary measure and \(1 < p \le 2\), so \(2 \le p^{\prime} < \infty\). This finiteness of \(p^{\prime}\) is essential below.

Local representations. For \(E \in \mathcal{S}\) with \(\mu(E) < \infty\), let \(\mu_E\) be the restriction of \(\mu\) to \(\mathcal{S}_E = \{ A \in \mathcal{S} : A \subseteq E \}\); this is a finite measure on \(E\). Extension by 0 off \(E\) is a linear isometry of \(L^p(\mu_E)\) into \(L^p(\mu)\); write \(\widetilde{f}\) for the extension of \(f\). Then \(f \mapsto \varphi(\widetilde{f})\) is a bounded linear functional on \(L^p(\mu_E)\) of norm at most \(\|\varphi\|\), so by part (a) there exists \(h_E \in L^{p^{\prime}}(\mu_E)\) with

\begin{equation*} \varphi(\widetilde{f}) = \int_E f h_E \, d\mu \qquad \text{for all } f \in L^p(\mu_E) . \end{equation*}

Because \(1 < p \le 2\) and \(\mu_E\) is a measure, 7.25 applies to \(\mu_E\): the map \(k \mapsto \varphi_k\) from \(L^{p^{\prime}}(\mu_E)\) to \(L^p(\mu_E)^{\prime}\) is one-to-one and norm-preserving. Consequently \(h_E\) is unique (as an element of \(L^{p^{\prime}}(\mu_E)\)) and

\begin{equation*} \|h_E\|_{p^{\prime}} = \| f \mapsto \varphi(\widetilde{f}) \, \| \le \|\varphi\| . \end{equation*}

Moreover, if \(E \subseteq F\) with \(\mu(F) < \infty\), then for every \(f \in L^p(\mu_E)\) the function \(\widetilde{f}\) restricted to \(F\) lies in \(L^p(\mu_F)\), so \(\int_E f \, h_F \, d\mu = \varphi(\widetilde{f}) = \int_E f h_E \, d\mu\); uniqueness gives

\begin{equation*} h_F |_E = h_E \quad \text{almost everywhere on } E . \end{equation*}

A maximizing \(\sigma\)-finite set. Let

\begin{equation*} M = \sup \{ \|h_E\|_{p^{\prime}} : E \in \mathcal{S}, \ \mu(E) < \infty \} \le \|\varphi\| < \infty . \end{equation*}

Choose \(F_n \in \mathcal{S}\) with \(\mu(F_n) < \infty\) and \(\|h_{F_n}\|_{p^{\prime}}^{p^{\prime}} > M^{p^{\prime}} - \frac{1}{n}\), and set \(E_n = F_1 \cup \cdots \cup F_n\). Then \(E_1 \subseteq E_2 \subseteq \cdots\), each \(\mu(E_n) < \infty\), and by the consistency relation \(\|h_{E_n}\|_{p^{\prime}}^{p^{\prime}} \ge \int_{F_n} |h_{F_n}|^{p^{\prime}} d\mu > M^{p^{\prime}} - \frac{1}{n}\); also \(\|h_{E_n}\|_{p^{\prime}} \le M\) and \(n \mapsto \|h_{E_n}\|_{p^{\prime}}\) is nondecreasing. Hence \(\|h_{E_n}\|_{p^{\prime}} \to M\). Let \(E = \bigcup_{n} E_n\), a \(\sigma\)-finite set.

Fix a representative of each \(h_{E_n}\) and let

\begin{equation*} N = \bigcup_{n < m} \{ x \in E_n : h_{E_m}(x) \neq h_{E_n}(x) \} . \end{equation*}

By the consistency relation each set in this countable union has measure 0, so \(\mu(N) = 0\). Replacing each representative \(h_{E_n}\) by \(h_{E_n} \chi_{X \setminus N}\) changes no element of \(L^{p^{\prime}}(\mu_{E_n})\) and no integral below, and after this replacement \(h_{E_m}(x) = h_{E_n}(x)\) for all \(n < m\) and all \(x \in E_n\). Hence the following prescription defines a function \(h\) on \(X\) unambiguously, and \(h\) is measurable because \(h = \lim_{n \to \infty} h_{E_n} \chi_{E_n}\) pointwise on \(X\):

\begin{equation*} h(x) = h_{E_n}(x) \text{ if } x \in E_n , \qquad h(x) = 0 \text{ if } x \notin E . \end{equation*}

The functions \(\chi_{E_n} |h|^{p^{\prime}}\) increase pointwise to \(|h|^{p^{\prime}}\), so by the Monotone Convergence Theorem (3.11),

\begin{equation*} \int_X |h|^{p^{\prime}} \, d\mu = \lim_{n \to \infty} \|h_{E_n}\|_{p^{\prime}}^{p^{\prime}} = M^{p^{\prime}} . \end{equation*}

Thus \(h \in L^{p^{\prime}}(\mu)\) with \(\|h\|_{p^{\prime}} = M \le \|\varphi\|\).

Nothing lives outside \(E\). Suppose \(A \in \mathcal{S}\), \(\mu(A) < \infty\), and \(A \cap E = \emptyset\). For each \(n\) the set \(E_n \cup A\) has finite measure, and \(h_{E_n \cup A}\) agrees with \(h_{E_n}\) on \(E_n\) and with \(h_A\) on \(A\); since \(E_n\) and \(A\) are disjoint,

\begin{equation*} M^{p^{\prime}} \ge \|h_{E_n \cup A}\|_{p^{\prime}}^{p^{\prime}} = \|h_{E_n}\|_{p^{\prime}}^{p^{\prime}} + \|h_A\|_{p^{\prime}}^{p^{\prime}} . \end{equation*}

Letting \(n \to \infty\) gives \(M^{p^{\prime}} \ge M^{p^{\prime}} + \|h_A\|_{p^{\prime}}^{p^{\prime}}\). Because \(M < \infty\) and \(p^{\prime} < \infty\), we conclude \(\|h_A\|_{p^{\prime}} = 0\), so \(h_A = 0\) and therefore

\begin{equation*} \varphi(\widetilde{f}) = 0 \qquad \text{for every } f \in L^p(\mu_A) . \end{equation*}

(Here \(p^{\prime} < \infty\) is essential: the additivity \(\|k\|_{p^{\prime}}^{p^{\prime}} = \|k \chi_{E_n}\|_{p^{\prime}}^{p^{\prime}} + \|k \chi_A\|_{p^{\prime}}^{p^{\prime}}\) over the disjoint pieces has no counterpart for \(p^{\prime} = \infty\), where the norm is a maximum rather than a sum.)

Conclusion. Let \(f \in L^p(\mu)\). First, \(\{ f \neq 0 \}\) is \(\sigma\)-finite: for \(k \in \mathbf{Z}^+\),

\begin{equation*} \tfrac{1}{k^p} \, \mu \bigl( \{ |f| > \tfrac{1}{k} \} \bigr) \le \int_{ \{ |f| > 1/k \} } |f|^p \, d\mu \le \|f\|_p^p < \infty , \end{equation*}

so each \(\{ |f| > \frac{1}{k} \}\) has finite measure, and \(\{ f \neq 0 \} = \bigcup_k \{ |f| > \frac{1}{k} \}\).

Write \(A = \{ f \neq 0 \} \setminus E\) and choose an increasing sequence \(A_1 \subseteq A_2 \subseteq \cdots\) of sets of finite measure with \(A = \bigcup_k A_k\); each \(A_k\) is disjoint from \(E\). By the previous paragraph \(\varphi(f \chi_{A_k}) = 0\) for every \(k\). Also

\begin{equation*} \|f \chi_A - f \chi_{A_k}\|_p^p = \int_{A \setminus A_k} |f|^p \, d\mu \longrightarrow 0 \end{equation*}

by the Dominated Convergence Theorem (3.31), so \(f \chi_{A_k} \to f \chi_A\) in \(L^p(\mu)\); continuity of \(\varphi\) gives \(\varphi(f \chi_A) = 0\). Since \(f\) vanishes off \(\{ f \neq 0 \}\) we have \(f \chi_{X \setminus E} = f \chi_A\), so

\begin{equation*} \varphi(f \chi_{X \setminus E}) = 0 = \int f \chi_{X \setminus E} \, h \, d\mu , \end{equation*}

the last equality because \(h = 0\) off \(E\).

Second, \(f \chi_{E_n}\) (viewed in \(L^p(\mu_{E_n})\) and extended by zero) satisfies

\begin{equation*} \varphi(f \chi_{E_n}) = \int_{E_n} f h_{E_n} \, d\mu = \int f \chi_{E_n} h \, d\mu . \end{equation*}

As above, \(\|f \chi_E - f \chi_{E_n}\|_p \to 0\) by the Dominated Convergence Theorem (3.31), so the left side converges to \(\varphi(f \chi_E)\) by continuity of \(\varphi\), while by Hölder’s inequality (7.9)

\begin{equation*} \Bigl| \int f \chi_E h \, d\mu - \int f \chi_{E_n} h \, d\mu \Bigr| \le \|f \chi_E - f \chi_{E_n}\|_p \, \|h\|_{p^{\prime}} \longrightarrow 0 . \end{equation*}

Hence \(\varphi(f \chi_E) = \int f \chi_E h \, d\mu\).

Adding the two displayed conclusions and using \(f = f \chi_E + f \chi_{X \setminus E}\) together with \(h = h \chi_E\):

\begin{equation*} \varphi(f) = \varphi(f \chi_E) + \varphi(f \chi_{X \setminus E}) = \int f \chi_E h \, d\mu = \int f h \, d\mu . \end{equation*}

Thus \(h \in L^{p^{\prime}}(\mu)\) represents \(\varphi\), completing the proof of (b).

Together with 7.25 this shows that for \(1 < p \le 2\) the map \(h \mapsto \varphi_h\) is an isometric isomorphism of \(L^{p^{\prime}}(\mu)\) onto \(L^p(\mu)^{\prime}\), so the dual of \(L^p(\mu)\) may be identified with \(L^{p^{\prime}}(\mu)\).

Problem (8B.26)

Prove that if \(V\) is an infinite-dimensional Hilbert space, then the Banach space \(\mathcal{B}(V, V)\) is nonseparable.

Solution

Recall that a metric space is called separable if it has a countable dense subset. We produce an uncountable family of elements of \(\mathcal{B}(V,V)\) whose pairwise distances are all at least 1; no metric space containing such a family can be separable.

Step 1: \(V\) contains an orthonormal sequence. Because \(V\) is infinite-dimensional, no finite list of elements of \(V\) spans \(V\). Hence we may choose inductively \(f_1, f_2, \dots \in V\) with

\begin{equation*} f_n \notin \operatorname{span}(f_1, \dots, f_{n-1}) \qquad \text{for each } n \in \mathbf{Z}^+ \end{equation*}

(for \(n = 1\) this just means \(f_1 \neq 0\)). Apply the Gram–Schmidt procedure: set \(e_1 = f_1 / \|f_1\|\) and, having defined orthonormal \(e_1, \dots, e_{n-1}\) with \(\operatorname{span}(e_1, \dots, e_{n-1}) = \operatorname{span}(f_1, \dots, f_{n-1})\), put

\begin{equation*} v_n = f_n - \sum_{k=1}^{n-1} \langle f_n , e_k \rangle e_k , \qquad e_n = \frac{v_n}{\|v_n\|} . \end{equation*}

Here \(v_n \neq 0\) because \(f_n \notin \operatorname{span}(f_1, \dots, f_{n-1}) = \operatorname{span}(e_1, \dots, e_{n-1})\). For \(m < n\) we have

\begin{equation*} \langle v_n , e_m \rangle = \langle f_n , e_m \rangle - \sum_{k=1}^{n-1} \langle f_n , e_k \rangle \langle e_k , e_m \rangle = \langle f_n , e_m \rangle - \langle f_n , e_m \rangle = 0 , \end{equation*}

so \(e_n\) is orthogonal to \(e_1, \dots, e_{n-1}\), and \(\|e_n\| = 1\). Also \(\operatorname{span}(e_1, \dots, e_n) = \operatorname{span}(f_1, \dots, f_n)\), so the induction continues. Thus \(e_1, e_2, \dots\) is an orthonormal sequence in \(V\).

Step 2: an uncountable family of orthogonal projections. For each subset \(S \subseteq \mathbf{Z}^+\) let

\begin{equation*} U_S = \overline{\operatorname{span} \{ e_k : k \in S \}} , \end{equation*}

a closed subspace of \(V\), and let \(T_S = P_{U_S} \in \mathcal{B}(V,V)\) be the orthogonal projection onto \(U_S\) (which exists because \(V\) is a Hilbert space and \(U_S\) is a closed subspace; \(T_S\) is linear by 8.37(c) and satisfies \(\|T_S f\| \le \|f\|\) by 8.37(d), so \(T_S\) is indeed a bounded linear map with \(\|T_S\| \le 1\)). Note \(P_{\{0\}} = 0\) handles \(S = \emptyset\).

We compute \(T_S e_k\) for each \(k \in \mathbf{Z}^+\).

If \(k \in S\), then \(e_k \in U_S\), so \(T_S e_k = e_k\) because \(P_U g = g\) for every closed subspace \(U\) and every \(g \in U\): taking \(f = g\) and \(h = g\) in 8.37(b), the element \(f - h = 0\) is orthogonal to every element of \(U\), so \(g = P_U g\).

If \(k \notin S\), let \(W_S = \operatorname{span}\{ e_j : j \in S \}\). Every \(g \in W_S\) has the form \(g = \sum_{j \in F} \alpha_j e_j\) for some finite \(F \subseteq S\), whence

\begin{equation*} \langle g , e_k \rangle = \sum_{j \in F} \alpha_j \langle e_j , e_k \rangle = 0 \end{equation*}

by orthonormality (note \(k \notin F\)). Thus \(e_k \in W_S^{\perp}\), and \(W_S^{\perp} = \overline{W_S}^{\, \perp} = U_S^{\perp}\) by 8.40(d). Since \(\operatorname{null} P_{U_S} = U_S^{\perp}\) by 8.45(a), we get \(T_S e_k = 0\).

Step 3: the family is \(1\)-separated. Suppose \(S, S^{\prime}\) are distinct subsets of \(\mathbf{Z}^+\). Choose \(k\) in the symmetric difference of \(S\) and \(S^{\prime}\); without loss of generality \(k \in S\) and \(k \notin S^{\prime}\). Then by Step 2,

\begin{equation*} (T_S - T_{S^{\prime}}) e_k = e_k - 0 = e_k , \end{equation*}

so, since \(\|e_k\| = 1\),

\begin{equation*} \|T_S - T_{S^{\prime}}\| \ge \|(T_S - T_{S^{\prime}}) e_k\| = 1 . \end{equation*}

Step 4: conclusion. The collection of all subsets of \(\mathbf{Z}^+\) is uncountable, by Cantor’s diagonal argument: if \(S_1, S_2, \dots\) were a list of all subsets of \(\mathbf{Z}^+\), then the set \(S = \{ k \in \mathbf{Z}^+ : k \notin S_k \}\) would satisfy \(S = S_m\) for some \(m \in \mathbf{Z}^+\), and then \(m \in S_m\) if and only if \(m \notin S_m\), a contradiction. By Step 3 the map \(S \mapsto T_S\) is one-to-one, so \(\{ T_S : S \subseteq \mathbf{Z}^+ \}\) is an uncountable subset of \(\mathcal{B}(V,V)\) any two distinct elements of which are at distance at least 1.

Now suppose, for contradiction, that \(D\) is a countable dense subset of \(\mathcal{B}(V,V)\). For each \(S \subseteq \mathbf{Z}^+\) choose \(d_S \in D\) with \(\|T_S - d_S\| < \frac{1}{2}\). If \(S \neq S^{\prime}\) and \(d_S = d_{S^{\prime}}\), then

\begin{equation*} \|T_S - T_{S^{\prime}}\| \le \|T_S - d_S\| + \|d_{S^{\prime}} - T_{S^{\prime}}\| < \tfrac{1}{2} + \tfrac{1}{2} = 1 , \end{equation*}

contradicting Step 3. Hence \(S \mapsto d_S\) is a one-to-one map from the uncountable collection of subsets of \(\mathbf{Z}^+\) into the countable set \(D\), which is impossible.

Therefore \(\mathcal{B}(V,V)\) has no countable dense subset, that is, \(\mathcal{B}(V,V)\) is nonseparable.

Exercises 8C

Problem (8C.1)

Verify that the family \(\{e_k\}_{k \in \mathbb{Z}}\) as defined in the third bullet point of Example 8.51 is an orthonormal family in \(L^2\big((-\pi, \pi]\big)\). The following formulas should help:

\begin{equation*} (\sin x)(\cos y) = \frac{\sin(x-y) + \sin(x+y)}{2}, \end{equation*}

\begin{equation*} (\sin x)(\sin y) = \frac{\cos(x-y) - \cos(x+y)}{2}, \end{equation*}

\begin{equation*} (\cos x)(\cos y) = \frac{\cos(x-y) + \cos(x+y)}{2}. \end{equation*}

Solution

Recall the definition from the third bullet point of 8.51: for \(k \in \mathbb{Z}\) the function \(e_k : (-\pi, \pi] \to \mathbb{R}\) is given by

\begin{equation*} e_k(t) = \begin{cases} \dfrac{1}{\sqrt{\pi}}\, \sin(kt) & \text{if } k > 0, \\[2mm] \dfrac{1}{\sqrt{2\pi}} & \text{if } k = 0, \\[2mm] \dfrac{1}{\sqrt{\pi}}\, \cos(kt) & \text{if } k < 0. \end{cases} \end{equation*}

Each \(e_k\) is continuous and bounded on the bounded interval \((-\pi, \pi]\), hence lies in \(L^2\big((-\pi, \pi]\big)\) with respect to Lebesgue measure. Because all these functions are real valued, the inner product is

\begin{equation*} \langle f, g \rangle = \int_{-\pi}^{\pi} f g \, d\lambda . \end{equation*}

Since \(\cos\) is an even function we have \(\cos(kt) = \cos(|k|t)\), so for \(k < 0\) we may write \(e_k(t) = \frac{1}{\sqrt{\pi}}\cos(|k|t)\). Thus the family consists of the constant \(\frac{1}{\sqrt{2\pi}}\) together with the functions \(\frac{1}{\sqrt{\pi}}\sin(mt)\) and \(\frac{1}{\sqrt{\pi}}\cos(mt)\) for \(m \in \mathbb{Z}^+\).

Basic integrals. For \(n \in \mathbb{Z}\) with \(n \neq 0\),

\begin{equation*} \int_{-\pi}^{\pi} \cos(nt)\, dt = \Big[\tfrac{1}{n}\sin(nt)\Big]_{-\pi}^{\pi} = \frac{\sin(n\pi) - \sin(-n\pi)}{n} = 0 \end{equation*}

and

\begin{equation*} \int_{-\pi}^{\pi} \sin(nt)\, dt = \Big[-\tfrac{1}{n}\cos(nt)\Big]_{-\pi}^{\pi} = \frac{\cos(-n\pi) - \cos(n\pi)}{n} = 0, \end{equation*}

the last equality because \(\cos\) is even. For \(n = 0\) we have \(\int_{-\pi}^{\pi}\cos(0 \cdot t)\, dt = 2\pi\) and \(\int_{-\pi}^{\pi}\sin(0 \cdot t)\, dt = 0\). In short: the integral over \((-\pi, \pi]\) of \(t \mapsto \sin(nt)\) is \(0\) for every \(n \in \mathbb{Z}\), and the integral of \(t \mapsto \cos(nt)\) is \(0\) unless \(n = 0\), in which case it is \(2\pi\).

Norm of \(e_0\). We have \(\|e_0\|^2 = \int_{-\pi}^{\pi} \frac{1}{2\pi}\, dt = 1\), so \(\|e_0\| = 1\).

Sine against sine. Let \(m, n \in \mathbb{Z}^+\). Taking \(x = mt\) and \(y = nt\) in the second displayed formula,

\begin{equation*} \int_{-\pi}^{\pi} \sin(mt)\sin(nt)\, dt = \frac{1}{2}\int_{-\pi}^{\pi} \cos\big((m-n)t\big)\, dt - \frac{1}{2}\int_{-\pi}^{\pi} \cos\big((m+n)t\big)\, dt . \end{equation*}

Since \(m + n \geq 2 \neq 0\), the second integral is \(0\). The first integral is \(0\) if \(m \neq n\) and equals \(2\pi\) if \(m = n\). Hence

\begin{equation*} \int_{-\pi}^{\pi} \sin(mt)\sin(nt)\, dt = \begin{cases} \pi & \text{if } m = n, \\ 0 & \text{if } m \neq n. \end{cases} \end{equation*}

Dividing by \(\pi\) shows that \(\langle e_m, e_n \rangle = 0\) for distinct \(m, n \in \mathbb{Z}^+\) and \(\|e_m\| = 1\) for \(m \in \mathbb{Z}^+\).

Cosine against cosine. Let \(m, n \in \mathbb{Z}^+\). Taking \(x = mt\), \(y = nt\) in the third displayed formula,

\begin{equation*} \int_{-\pi}^{\pi} \cos(mt)\cos(nt)\, dt = \frac{1}{2}\int_{-\pi}^{\pi} \cos\big((m-n)t\big)\, dt + \frac{1}{2}\int_{-\pi}^{\pi} \cos\big((m+n)t\big)\, dt, \end{equation*}

which by the same reasoning equals \(\pi\) if \(m = n\) and \(0\) if \(m \neq n\). Dividing by \(\pi\) shows that \(\langle e_{-m}, e_{-n}\rangle = 0\) for distinct \(m, n \in \mathbb{Z}^+\) and \(\|e_{-m}\| = 1\) for \(m \in \mathbb{Z}^+\).

Sine against cosine. Let \(m, n \in \mathbb{Z}^+\). The first displayed formula with \(x = mt\), \(y = nt\) gives

\begin{equation*} \int_{-\pi}^{\pi} \sin(mt)\cos(nt)\, dt = \frac{1}{2}\int_{-\pi}^{\pi} \sin\big((m-n)t\big)\, dt + \frac{1}{2}\int_{-\pi}^{\pi} \sin\big((m+n)t\big)\, dt = 0, \end{equation*}

because the integral of \(t \mapsto \sin(nt)\) over \((-\pi, \pi]\) vanishes for every integer \(n\) (including \(n = 0\)). Hence \(\langle e_m, e_{-n} \rangle = 0\) for all \(m, n \in \mathbb{Z}^+\).

Everything against \(e_0\). For \(m \in \mathbb{Z}^+\),

\begin{equation*} \langle e_m, e_0 \rangle = \frac{1}{\sqrt{2}\,\pi}\int_{-\pi}^{\pi} \sin(mt)\, dt = 0, \qquad \langle e_{-m}, e_0 \rangle = \frac{1}{\sqrt{2}\,\pi}\int_{-\pi}^{\pi} \cos(mt)\, dt = 0 . \end{equation*}

The five computations above cover every pair \(j, k \in \mathbb{Z}\): we have shown \(\|e_k\| = 1\) for all \(k \in \mathbb{Z}\) and \(\langle e_j, e_k\rangle = 0\) whenever \(j \neq k\). Thus \(\{e_k\}_{k \in \mathbb{Z}}\) satisfies the definition 8.50 of an orthonormal family in \(L^2\big((-\pi, \pi]\big)\).

Problem (8C.2)

Suppose \(\{a_k\}_{k \in \Gamma}\) is a family in \(\mathbb{R}\) and \(a_k \geq 0\) for each \(k \in \Gamma\). Prove the unordered sum \(\sum_{k \in \Gamma} a_k\) converges if and only if

\begin{equation*} \sup\Big\{ \sum_{j \in \Omega} a_j : \Omega \text{ is a finite subset of } \Gamma \Big\} < \infty . \end{equation*}

Furthermore, prove that if \(\sum_{k \in \Gamma} a_k\) converges then it equals the supremum above.

Solution

Write

\begin{equation*} S = \sup\Big\{ \sum_{j \in \Omega} a_j : \Omega \text{ is a finite subset of } \Gamma \Big\} \in [0, \infty], \end{equation*}

where the supremum is over all finite subsets \(\Omega\) of \(\Gamma\) (taking \(\Omega = \varnothing\) gives the value \(0\), so \(S \geq 0\)).

The key elementary observation is monotonicity: if \(\Omega \subseteq \Omega^{\prime}\) are finite subsets of \(\Gamma\), then

\begin{equation*} \sum_{j \in \Omega^{\prime}} a_j - \sum_{j \in \Omega} a_j = \sum_{j \in \Omega^{\prime} \setminus \Omega} a_j \geq 0, \end{equation*}

because every \(a_j\) is nonnegative. Thus \(\sum_{j \in \Omega} a_j \leq \sum_{j \in \Omega^{\prime}} a_j\).

Suppose first that \(S < \infty\). We show \(\sum_{k \in \Gamma} a_k\) converges, with value \(S\). Let \(\varepsilon > 0\). By the definition of supremum there is a finite subset \(\Omega\) of \(\Gamma\) with

\begin{equation*} \sum_{j \in \Omega} a_j > S - \varepsilon . \end{equation*}

Now let \(\Omega^{\prime}\) be any finite set with \(\Omega \subseteq \Omega^{\prime} \subseteq \Gamma\). By monotonicity and by the definition of \(S\),

\begin{equation*} S - \varepsilon < \sum_{j \in \Omega} a_j \leq \sum_{j \in \Omega^{\prime}} a_j \leq S . \end{equation*}

Hence \(\big| S - \sum_{j \in \Omega^{\prime}} a_j \big| < \varepsilon\) for every finite \(\Omega^{\prime}\) with \(\Omega \subseteq \Omega^{\prime} \subseteq \Gamma\). Since \(\varepsilon > 0\) was arbitrary, the definition 8.53 of an unordered sum shows that \(\sum_{k \in \Gamma} a_k\) converges and equals \(S\).

Conversely, suppose \(\sum_{k \in \Gamma} a_k\) converges, say \(\sum_{k \in \Gamma} a_k = g \in \mathbb{R}\). Let \(\varepsilon > 0\). By 8.53 there exists a finite subset \(\Omega_\varepsilon\) of \(\Gamma\) such that

\begin{equation*} \Big| g - \sum_{j \in \Omega^{\prime}} a_j \Big| < \varepsilon \end{equation*}

for all finite sets \(\Omega^{\prime}\) with \(\Omega_\varepsilon \subseteq \Omega^{\prime} \subseteq \Gamma\).

Let \(\Omega\) be an arbitrary finite subset of \(\Gamma\). Then \(\Omega \cup \Omega_\varepsilon\) is a finite subset of \(\Gamma\) containing \(\Omega_\varepsilon\), so the displayed inequality applies to it; combining with monotonicity gives

\begin{equation*} \sum_{j \in \Omega} a_j \leq \sum_{j \in \Omega \cup \Omega_\varepsilon} a_j < g + \varepsilon . \end{equation*}

Because \(\Omega\) was an arbitrary finite subset of \(\Gamma\), this shows \(S \leq g + \varepsilon\). In particular, taking \(\varepsilon = 1\) gives \(S \leq g + 1 < \infty\), which proves that the supremum is finite.

Finally we show \(S = g\). On the one hand, the previous paragraph gives \(S \leq g + \varepsilon\) for every \(\varepsilon > 0\), hence \(S \leq g\). On the other hand, applying the definition of the unordered sum to the finite set \(\Omega_\varepsilon\) itself gives

\begin{equation*} S \geq \sum_{j \in \Omega_\varepsilon} a_j > g - \varepsilon \end{equation*}

for every \(\varepsilon > 0\), hence \(S \geq g\). Therefore \(S = g\), which is exactly the assertion that if \(\sum_{k \in \Gamma} a_k\) converges then it equals the supremum \(S\).

Problem (8C.3)

Suppose \(\{e_k\}_{k \in \Gamma}\) is an orthonormal family in an inner product space \(V\). Prove that if \(f \in V\), then \(\{k \in \Gamma : \langle f, e_k \rangle \neq 0\}\) is a countable set.

Solution

Fix \(f \in V\).

First we record the finite form of Bessel’s inequality: if \(\Omega\) is a finite subset of \(\Gamma\), then

\begin{equation*} \sum_{j \in \Omega} |\langle f, e_j \rangle|^2 \leq \|f\|^2 . \end{equation*}

This is immediate from Bessel’s inequality 8.57 together with Exercise 2 of this section, which identifies the unordered sum of a nonnegative family with the supremum of its finite subsums. It also has a two-line direct proof: put \(g = \sum_{j \in \Omega} \langle f, e_j\rangle e_j\); for each \(k \in \Omega\) orthonormality gives \(\langle g, e_k \rangle = \langle f, e_k \rangle\), hence \(\langle f - g, e_k\rangle = 0\) for all \(k \in \Omega\) and therefore \(\langle f - g, g \rangle = 0\). The Pythagorean theorem and 8.52 then give

\begin{equation*} \|f\|^2 = \|f - g\|^2 + \|g\|^2 \geq \|g\|^2 = \sum_{j \in \Omega} |\langle f, e_j \rangle|^2 . \end{equation*}

Now for each \(n \in \mathbb{Z}^+\) define

\begin{equation*} \Gamma_n = \Big\{ k \in \Gamma : |\langle f, e_k \rangle|^2 > \tfrac{1}{n} \Big\} . \end{equation*}

We claim each \(\Gamma_n\) is a finite set with at most \(n\|f\|^2\) elements. To see this, suppose \(\Omega\) is a nonempty finite subset of \(\Gamma_n\), say with \(m \geq 1\) elements. Every \(j \in \Omega\) lies in \(\Gamma_n\) and so contributes more than \(1/n\) to the sum below; hence

\begin{equation*} \frac{m}{n} < \sum_{j \in \Omega} |\langle f, e_j \rangle|^2 \leq \|f\|^2, \end{equation*}

the second inequality being the finite Bessel inequality proved above. Thus \(m < n\|f\|^2\), and in particular \(m \leq n\|f\|^2\). Since the empty subset trivially has at most \(n\|f\|^2\) elements, every finite subset of \(\Gamma_n\) has at most \(n\|f\|^2\) elements.

This forces \(\Gamma_n\) itself to be finite with at most \(n\|f\|^2\) elements. Indeed, if \(\Gamma_n\) had more than \(n\|f\|^2\) elements (in particular, if \(\Gamma_n\) were infinite), we could choose finitely many distinct elements of \(\Gamma_n\), more than \(n\|f\|^2\) of them, and the resulting finite subset would contradict the previous paragraph.

Finally, if \(\langle f, e_k \rangle \neq 0\) then \(|\langle f, e_k\rangle|^2 > 0\), so \(|\langle f, e_k\rangle|^2 > 1/n\) for some \(n \in \mathbb{Z}^+\) by the Archimedean property; conversely each \(\Gamma_n\) consists of indices with \(\langle f, e_k\rangle \neq 0\). Therefore

\begin{equation*} \{k \in \Gamma : \langle f, e_k \rangle \neq 0\} = \bigcup_{n = 1}^{\infty} \Gamma_n . \end{equation*}

This exhibits \(\{k \in \Gamma : \langle f, e_k\rangle \neq 0\}\) as a countable union of finite sets, and a countable union of countable sets is countable. Hence \(\{k \in \Gamma : \langle f, e_k \rangle \neq 0\}\) is countable, as desired.

Problem (8C.4)

Suppose \(\{f_k\}_{k \in \Gamma}\) and \(\{g_k\}_{k \in \Gamma}\) are families in a normed vector space such that \(\sum_{k \in \Gamma} f_k\) and \(\sum_{k \in \Gamma} g_k\) converge. Prove that \(\sum_{k \in \Gamma}(f_k + g_k)\) converges and

\begin{equation*} \sum_{k \in \Gamma} (f_k + g_k) = \sum_{k \in \Gamma} f_k + \sum_{k \in \Gamma} g_k . \end{equation*}

Solution

Let \(V\) denote the normed vector space in question, and set

\begin{equation*} f = \sum_{k \in \Gamma} f_k, \qquad g = \sum_{k \in \Gamma} g_k, \end{equation*}

both of which exist in \(V\) by hypothesis. We must show that the unordered sum \(\sum_{k \in \Gamma}(f_k + g_k)\) converges to \(f + g\), in the sense of definition 8.53.

Suppose \(\varepsilon > 0\). Because \(\sum_{k \in \Gamma} f_k = f\), there is a finite subset \(\Omega_1\) of \(\Gamma\) such that

\begin{equation*} \Big\| f - \sum_{j \in \Omega^{\prime}} f_j \Big\| < \frac{\varepsilon}{2} \end{equation*}

for every finite set \(\Omega^{\prime}\) with \(\Omega_1 \subseteq \Omega^{\prime} \subseteq \Gamma\). Because \(\sum_{k \in \Gamma} g_k = g\), there is likewise a finite subset \(\Omega_2\) of \(\Gamma\) such that

\begin{equation*} \Big\| g - \sum_{j \in \Omega^{\prime}} g_j \Big\| < \frac{\varepsilon}{2} \end{equation*}

for every finite set \(\Omega^{\prime}\) with \(\Omega_2 \subseteq \Omega^{\prime} \subseteq \Gamma\).

Let \(\Omega = \Omega_1 \cup \Omega_2\), which is a finite subset of \(\Gamma\). Suppose \(\Omega^{\prime}\) is a finite set with \(\Omega \subseteq \Omega^{\prime} \subseteq \Gamma\). Then \(\Omega_1 \subseteq \Omega^{\prime}\) and \(\Omega_2 \subseteq \Omega^{\prime}\), so both displayed inequalities apply to \(\Omega^{\prime}\). Because \(\Omega^{\prime}\) is finite, the sum \(\sum_{j \in \Omega^{\prime}}(f_j + g_j)\) is an ordinary finite sum in a vector space, and commutativity and associativity of vector addition give

\begin{equation*} \sum_{j \in \Omega^{\prime}} (f_j + g_j) = \sum_{j \in \Omega^{\prime}} f_j + \sum_{j \in \Omega^{\prime}} g_j . \end{equation*}

Hence, by the triangle inequality,

\begin{equation*} \begin{aligned} \Big\| (f + g) - \sum_{j \in \Omega^{\prime}} (f_j + g_j) \Big\| &= \Big\| \Big(f - \sum_{j \in \Omega^{\prime}} f_j\Big) + \Big(g - \sum_{j \in \Omega^{\prime}} g_j\Big) \Big\| \\ &\leq \Big\| f - \sum_{j \in \Omega^{\prime}} f_j \Big\| + \Big\| g - \sum_{j \in \Omega^{\prime}} g_j \Big\| \\ &< \frac{\varepsilon}{2} + \frac{\varepsilon}{2} = \varepsilon . \end{aligned} \end{equation*}

Thus for each \(\varepsilon > 0\) we have produced a finite subset \(\Omega\) of \(\Gamma\) such that \(\big\|(f+g) - \sum_{j \in \Omega^{\prime}}(f_j + g_j)\big\| < \varepsilon\) for all finite \(\Omega^{\prime}\) with \(\Omega \subseteq \Omega^{\prime} \subseteq \Gamma\). By 8.53 this says exactly that \(\sum_{k \in \Gamma}(f_k + g_k)\) converges and

\begin{equation*} \sum_{k \in \Gamma} (f_k + g_k) = f + g = \sum_{k \in \Gamma} f_k + \sum_{k \in \Gamma} g_k, \end{equation*}

as desired.

Problem (8C.5)

Suppose \(\{f_k\}_{k \in \Gamma}\) is a family in a normed vector space such that \(\sum_{k \in \Gamma} f_k\) converges. Prove that if \(c \in \mathbf{F}\), then \(\sum_{k \in \Gamma}(c f_k)\) converges and

\begin{equation*} \sum_{k \in \Gamma} (c f_k) = c \sum_{k \in \Gamma} f_k . \end{equation*}

Solution

Let \(f = \sum_{k \in \Gamma} f_k\), which exists in the normed vector space \(V\) by hypothesis, and let \(c \in \mathbf{F}\). We show that \(\sum_{k \in \Gamma} (c f_k)\) converges to \(c f\) in the sense of 8.53.

Suppose \(\varepsilon > 0\). Because \(\sum_{k \in \Gamma} f_k = f\), there is a finite subset \(\Omega\) of \(\Gamma\) such that

\begin{equation*} \Big\| f - \sum_{j \in \Omega^{\prime}} f_j \Big\| < \frac{\varepsilon}{1 + |c|} \end{equation*}

for every finite set \(\Omega^{\prime}\) with \(\Omega \subseteq \Omega^{\prime} \subseteq \Gamma\) (note \(\varepsilon/(1+|c|) > 0\), so this is a legitimate application of the definition of the unordered sum).

Let \(\Omega^{\prime}\) be a finite set with \(\Omega \subseteq \Omega^{\prime} \subseteq \Gamma\). Since \(\Omega^{\prime}\) is finite, distributivity of scalar multiplication over finite sums gives

\begin{equation*} \sum_{j \in \Omega^{\prime}} (c f_j) = c \sum_{j \in \Omega^{\prime}} f_j . \end{equation*}

Therefore, using homogeneity of the norm,

\begin{equation*} \Big\| cf - \sum_{j \in \Omega^{\prime}} (c f_j) \Big\| = \Big\| c \Big( f - \sum_{j \in \Omega^{\prime}} f_j \Big) \Big\| = |c| \, \Big\| f - \sum_{j \in \Omega^{\prime}} f_j \Big\| \leq \frac{|c|}{1 + |c|} \, \varepsilon < \varepsilon . \end{equation*}

Thus for every \(\varepsilon > 0\) there is a finite subset \(\Omega\) of \(\Gamma\) such that \(\big\| cf - \sum_{j \in \Omega^{\prime}}(cf_j)\big\| < \varepsilon\) for all finite \(\Omega^{\prime}\) with \(\Omega \subseteq \Omega^{\prime} \subseteq \Gamma\). By 8.53, the unordered sum \(\sum_{k \in \Gamma}(c f_k)\) converges and equals \(cf\), which is the desired conclusion

\begin{equation*} \sum_{k \in \Gamma} (c f_k) = c \sum_{k \in \Gamma} f_k . \end{equation*}

Problem (8C.6)

Suppose \(\{a_k\}_{k \in \Gamma}\) is a family in \(\mathbb{R}\). Prove that the unordered sum \(\sum_{k \in \Gamma} a_k\) converges if and only if \(\sum_{k \in \Gamma} |a_k| < \infty\).

Solution

Throughout, we use Exercise 2 of this section: for a family of nonnegative real numbers, the unordered sum converges if and only if the supremum of its finite subsums is finite, in which case the unordered sum equals that supremum. As in the text, for a nonnegative family the statement \(\sum_{k \in \Gamma} |a_k| < \infty\) means precisely that this unordered sum converges, equivalently that

\begin{equation*} S = \sup\Big\{ \sum_{j \in \Omega} |a_j| : \Omega \text{ is a finite subset of } \Gamma \Big\} < \infty . \end{equation*}

Suppose first that \(S < \infty\). For each \(k \in \Gamma\) put

\begin{equation*} a_k^+ = \max\{a_k, 0\}, \qquad a_k^- = \max\{-a_k, 0\}, \end{equation*}

so that \(a_k^+ \geq 0\), \(a_k^- \geq 0\), \(a_k = a_k^+ - a_k^-\), and \(a_k^+ \le |a_k|\), \(a_k^- \le |a_k|\). If \(\Omega\) is any finite subset of \(\Gamma\), then

\begin{equation*} \sum_{j \in \Omega} a_j^+ \leq \sum_{j \in \Omega} |a_j| \leq S \qquad\text{and}\qquad \sum_{j \in \Omega} a_j^- \leq \sum_{j \in \Omega} |a_j| \leq S . \end{equation*}

Thus the suprema of the finite subsums of the nonnegative families \(\{a_k^+\}_{k \in \Gamma}\) and \(\{a_k^-\}_{k \in \Gamma}\) are both at most \(S\), hence finite, so by Exercise 2 both unordered sums \(\sum_{k \in \Gamma} a_k^+\) and \(\sum_{k \in \Gamma} a_k^-\) converge. By Exercise 5 (with \(c = -1\)), \(\sum_{k \in \Gamma}(-a_k^-)\) converges, and then by Exercise 4 the unordered sum

\begin{equation*} \sum_{k \in \Gamma} a_k = \sum_{k \in \Gamma} \big( a_k^+ + (-a_k^-) \big) \end{equation*}

converges (and equals \(\sum_{k \in \Gamma} a_k^+ - \sum_{k \in \Gamma} a_k^-\)).

Conversely, suppose \(\sum_{k \in \Gamma} a_k\) converges, say to \(g \in \mathbb{R}\). Applying 8.53 with \(\varepsilon = 1\), there is a finite subset \(\Omega_0\) of \(\Gamma\) such that

\begin{equation*} \Big| g - \sum_{j \in \Omega^{\prime}} a_j \Big| < 1 \end{equation*}

for every finite set \(\Omega^{\prime}\) with \(\Omega_0 \subseteq \Omega^{\prime} \subseteq \Gamma\).

Claim: if \(\Omega\) is a finite subset of \(\Gamma \setminus \Omega_0\), then \(\big|\sum_{j \in \Omega} a_j\big| < 2\). Indeed, \(\Omega_0\) and \(\Omega_0 \cup \Omega\) are both finite sets containing \(\Omega_0\) and contained in \(\Gamma\), and since \(\Omega \cap \Omega_0 = \varnothing\) we have

\begin{equation*} \sum_{j \in \Omega} a_j = \sum_{j \in \Omega_0 \cup \Omega} a_j - \sum_{j \in \Omega_0} a_j = \Big( \sum_{j \in \Omega_0 \cup \Omega} a_j - g \Big) + \Big( g - \sum_{j \in \Omega_0} a_j \Big), \end{equation*}

whose absolute value is less than \(1 + 1 = 2\) by the displayed inequality.

Now let \(\Omega\) be an arbitrary finite subset of \(\Gamma \setminus \Omega_0\), and split it as \(\Omega = \Omega^+ \cup \Omega^-\), where \(\Omega^+ = \{j \in \Omega : a_j \geq 0\}\) and \(\Omega^- = \{j \in \Omega : a_j < 0\}\). Both \(\Omega^+\) and \(\Omega^-\) are finite subsets of \(\Gamma \setminus \Omega_0\), so the claim applies to each, giving

\begin{equation*} \sum_{j \in \Omega} |a_j| = \sum_{j \in \Omega^+} a_j + \sum_{j \in \Omega^-} (-a_j) = \Big| \sum_{j \in \Omega^+} a_j \Big| + \Big| \sum_{j \in \Omega^-} a_j \Big| < 2 + 2 = 4 . \end{equation*}

Finally, let \(\Omega\) be an arbitrary finite subset of \(\Gamma\). Then \(\Omega \setminus \Omega_0\) is a finite subset of \(\Gamma \setminus \Omega_0\), so

\begin{equation*} \sum_{j \in \Omega} |a_j| \leq \sum_{j \in \Omega_0} |a_j| + \sum_{j \in \Omega \setminus \Omega_0} |a_j| < \sum_{j \in \Omega_0} |a_j| + 4 . \end{equation*}

The right side is a fixed finite number, because \(\Omega_0\) is finite. Hence \(S \leq \sum_{j \in \Omega_0}|a_j| + 4 < \infty\), and by Exercise 2 applied to the nonnegative family \(\{|a_k|\}_{k \in \Gamma}\) we conclude that \(\sum_{k \in \Gamma} |a_k|\) converges with \(\sum_{k \in \Gamma}|a_k| = S < \infty\).

Thus convergence of an unordered sum of real numbers is equivalent to absolute convergence, as asserted in the text following 8.53.

Problem (8C.7)

Suppose \(\{f_k\}_{k \in \mathbb{Z}^+}\) is a family in a normed vector space \(V\) and \(f \in V\). Prove that the unordered sum \(\sum_{k \in \mathbb{Z}^+} f_k\) equals \(f\) if and only if the usual ordered sum \(\sum_{k=1}^{\infty} f_{p(k)}\) equals \(f\) for every injective and surjective function \(p : \mathbb{Z}^+ \to \mathbb{Z}^+\).

Solution

Recall (6.40) that the ordered sum \(\sum_{k=1}^{\infty} g_k\) equals \(f\) means \(\lim_{n \to \infty} \sum_{k=1}^{n} g_k = f\) in the norm of \(V\). Note also that if \(p\) is injective and \(n \in \mathbb{Z}^+\), then the finite sum \(\sum_{k=1}^{n} f_{p(k)}\) is exactly the finite unordered sum \(\sum_{j \in \Omega} f_j\) over the \(n\)-element set \(\Omega = \{p(1), \dots, p(n)\}\), since finite sums in a vector space do not depend on the order of the terms.

First suppose \(\sum_{k \in \mathbb{Z}^+} f_k = f\), and let \(p : \mathbb{Z}^+ \to \mathbb{Z}^+\) be a bijection. Let \(\varepsilon > 0\). By 8.53 there is a finite subset \(\Omega\) of \(\mathbb{Z}^+\) such that

\begin{equation*} \Big\| f - \sum_{j \in \Omega^{\prime}} f_j \Big\| < \varepsilon \end{equation*}

for all finite sets \(\Omega^{\prime}\) with \(\Omega \subseteq \Omega^{\prime} \subseteq \mathbb{Z}^+\). Because \(p\) is a bijection, \(p^{-1}(j)\) is defined for each \(j \in \Omega\); since \(\Omega\) is finite we may set

\begin{equation*} N = \max\{ p^{-1}(j) : j \in \Omega \} \end{equation*}

(with \(N = 1\) if \(\Omega = \varnothing\)). If \(n \geq N\), then every \(j \in \Omega\) satisfies \(j = p(p^{-1}(j))\) with \(p^{-1}(j) \leq N \leq n\), so \(\Omega \subseteq \{p(1), \dots, p(n)\}\). Applying the display to the finite set \(\Omega^{\prime} = \{p(1), \dots, p(n)\}\) gives

\begin{equation*} \Big\| f - \sum_{k=1}^{n} f_{p(k)} \Big\| = \Big\| f - \sum_{j \in \Omega^{\prime}} f_j \Big\| < \varepsilon \end{equation*}

for all \(n \geq N\). Since \(\varepsilon > 0\) was arbitrary, \(\lim_{n \to \infty} \sum_{k=1}^{n} f_{p(k)} = f\); that is, \(\sum_{k=1}^{\infty} f_{p(k)} = f\).

For the other direction we prove the contrapositive: assuming it is not true that \(\sum_{k \in \mathbb{Z}^+} f_k = f\), we construct a bijection \(p : \mathbb{Z}^+ \to \mathbb{Z}^+\) for which \(\sum_{k=1}^{\infty} f_{p(k)}\) does not equal \(f\).

Negating 8.53, there exists \(\varepsilon > 0\) with the following property: for every finite subset \(A\) of \(\mathbb{Z}^+\) there is a finite set \(\Omega\) with \(A \subseteq \Omega \subseteq \mathbb{Z}^+\) and

\begin{equation*} \Big\| f - \sum_{j \in \Omega} f_j \Big\| \geq \varepsilon . \end{equation*}

Using this, define finite sets \(\Omega_1 \subseteq \Omega_2 \subseteq \cdots\) inductively. Put \(\Omega_0 = \varnothing\). Given the finite set \(\Omega_{m-1}\), apply the property above to the finite set \(A_m = \Omega_{m-1} \cup \{1, 2, \dots, m\}\), obtaining a finite set \(\Omega_m\) with \(A_m \subseteq \Omega_m \subseteq \mathbb{Z}^+\) and \(\big\| f - \sum_{j \in \Omega_m} f_j \big\| \geq \varepsilon\). By construction \(\Omega_{m-1} \subseteq \Omega_m\) and \(\{1, \dots, m\} \subseteq \Omega_m\) for each \(m \in \mathbb{Z}^+\); hence the \(\Omega_m\) increase, \(\bigcup_{m=1}^{\infty} \Omega_m = \mathbb{Z}^+\), and \(n_m := |\Omega_m| \geq m\), so \(n_m \to \infty\).

Now define \(p\) by listing the elements of \(\mathbb{Z}^+\) block by block: let \(p(1), \dots, p(n_1)\) be the elements of \(\Omega_1\) in increasing order, and for each \(m \geq 2\) let \(p(n_{m-1}+1), \dots, p(n_m)\) be the elements of \(\Omega_m \setminus \Omega_{m-1}\) in increasing order (this block is empty if \(\Omega_m = \Omega_{m-1}\), in which case \(n_m = n_{m-1}\) and nothing is listed at this stage). Since \(n_m \to \infty\), this defines \(p(n)\) for every \(n \in \mathbb{Z}^+\). The blocks are pairwise disjoint and each element is listed once, so \(p\) is injective; every \(j \in \mathbb{Z}^+\) lies in some \(\Omega_m\) and hence is listed in one of the first \(m\) blocks, so \(p\) is surjective. Thus \(p : \mathbb{Z}^+ \to \mathbb{Z}^+\) is a bijection, and by construction

\begin{equation*} \{ p(1), \dots, p(n_m) \} = \Omega_m \quad\text{for every } m \in \mathbb{Z}^+ . \end{equation*}

Consequently, for every \(m \in \mathbb{Z}^+\),

\begin{equation*} \Big\| f - \sum_{k=1}^{n_m} f_{p(k)} \Big\| = \Big\| f - \sum_{j \in \Omega_m} f_j \Big\| \geq \varepsilon . \end{equation*}

Because \(n_1 \leq n_2 \leq \cdots\) and \(n_m \to \infty\), the sequence of partial sums \(\big(\sum_{k=1}^{n} f_{p(k)}\big)_{n \in \mathbb{Z}^+}\) has terms arbitrarily far out that stay at distance at least \(\varepsilon\) from \(f\). Hence this sequence does not converge to \(f\), so \(\sum_{k=1}^{\infty} f_{p(k)}\) does not equal \(f\).

This proves the contrapositive of the remaining direction: if \(\sum_{k=1}^{\infty} f_{p(k)} = f\) for every bijection \(p : \mathbb{Z}^+ \to \mathbb{Z}^+\), then \(\sum_{k \in \mathbb{Z}^+} f_k = f\). Combining the two directions completes the proof.

Problem (8C.8)

Explain why 8.58 implies that if \(\Gamma\) is a finite set and \(\{e_k\}_{k \in \Gamma}\) is an orthonormal family in a Hilbert space \(V\), then \(\operatorname{span}\{e_k\}_{k \in \Gamma}\) is a closed subspace of \(V\).

Solution

Recall (6.54) that the span of a family \(\{e_k\}_{k \in \Gamma}\) is the set of all sums \(\sum_{j \in \Omega} \alpha_j e_j\), where \(\Omega\) is a finite subset of \(\Gamma\) and \(\{\alpha_j\}_{j \in \Omega}\) is a family in \(\mathbf{F}\). In particular the span is a subspace of \(V\); the only issue is whether it is closed.

Suppose now that \(\Gamma\) is finite. Then given a finite \(\Omega \subseteq \Gamma\) and a family \(\{\alpha_j\}_{j \in \Omega}\) in \(\mathbf{F}\), we may extend the family to all of \(\Gamma\) by setting \(\alpha_j = 0\) for \(j \in \Gamma \setminus \Omega\), without changing the value of the sum. Hence

\begin{equation*} \operatorname{span}\{e_k\}_{k \in \Gamma} = \Big\{ \sum_{k \in \Gamma} \alpha_k e_k : \{\alpha_k\}_{k \in \Gamma} \text{ is a family in } \mathbf{F} \Big\}, \end{equation*}

where the sum on the right is an ordinary finite sum.

Next observe two things, both consequences of the finiteness of \(\Gamma\).

First, every family \(\{\alpha_k\}_{k \in \Gamma}\) in \(\mathbf{F}\) satisfies \(\sum_{k \in \Gamma} |\alpha_k|^2 < \infty\), because that sum has only finitely many terms, each finite.

Second, for a family \(\{\alpha_k\}_{k \in \Gamma}\) in \(\mathbf{F}\) the unordered sum \(\sum_{k \in \Gamma} \alpha_k e_k\) in the sense of 8.53 converges and equals the ordinary finite sum \(g = \sum_{k \in \Gamma} \alpha_k e_k\). Indeed, given \(\varepsilon > 0\) take the finite set in Definition 8.53 to be \(\Omega = \Gamma\) itself; the only finite set \(\Omega^{\prime}\) with \(\Gamma \subseteq \Omega^{\prime} \subseteq \Gamma\) is \(\Gamma\), and \(\|g - \sum_{j \in \Gamma} \alpha_j e_j\| = 0 < \varepsilon\).

Combining these two observations, the set appearing on the right side of 8.58(a), namely

\begin{equation*} \Big\{ \sum_{k \in \Gamma} \alpha_k e_k : \{\alpha_k\}_{k \in \Gamma} \text{ is a family in } \mathbf{F} \text{ and } \sum_{k \in \Gamma} |\alpha_k|^2 < \infty \Big\}, \end{equation*}

is exactly \(\operatorname{span}\{e_k\}_{k \in \Gamma}\): the constraint \(\sum_{k \in \Gamma} |\alpha_k|^2 < \infty\) is vacuous, and each unordered sum is an ordinary finite linear combination of \(\{e_k\}_{k \in \Gamma}\).

Therefore 8.58(a) reads

\begin{equation*} \overline{\operatorname{span}\{e_k\}_{k \in \Gamma}} = \operatorname{span}\{e_k\}_{k \in \Gamma}. \end{equation*}

A subset of a metric space equals its own closure if and only if it is closed (see 6.7 and the discussion following it). Hence \(\operatorname{span}\{e_k\}_{k \in \Gamma}\) is a closed subspace of \(V\), as desired.

Note that the finiteness of \(\Gamma\) was used twice and cannot be dropped: if \(\Gamma\) is infinite, then 8.58(a) exhibits elements of the closure of the span, such as \(\sum_{n} \frac{1}{n} e_{k_n}\) for distinct \(k_1, k_2, \ldots \in \Gamma\), that are not finite linear combinations of the \(e_k\) (see Exercise 9).

Problem (8C.9)

Suppose \(V\) is an infinite-dimensional Hilbert space. Prove that there does not exist a basis of \(V\) that is an orthonormal family.

Solution

Suppose, seeking a contradiction, that \(\{e_k\}_{k \in \Gamma}\) is an orthonormal family in \(V\) that is a basis of \(V\); by 6.54 this means the family is linearly independent and

\begin{equation*} \operatorname{span}\{e_k\}_{k \in \Gamma} = V. \end{equation*}

First we note that \(\Gamma\) must be infinite. Indeed, if \(\Gamma\) were finite then \(V\) would be the span of the finite list \(\{e_k\}_{k \in \Gamma}\) and hence would be finite-dimensional, contradicting our hypothesis.

Because \(\Gamma\) is infinite, we may choose a sequence \(k_1, k_2, \ldots\) of distinct elements of \(\Gamma\). Define a family \(\{\alpha_k\}_{k \in \Gamma}\) in \(\mathbf{F}\) by

\begin{equation*} \alpha_k = \begin{cases} \frac{1}{n} & \text{if } k = k_n \text{ for some } n \in \mathbf{Z}^+, \\ 0 & \text{otherwise.}\end{cases} \end{equation*}

This is well defined because \(k_1, k_2, \ldots\) are distinct. We have

\begin{equation*} \sum_{k \in \Gamma} |\alpha_k|^2 = \sum_{n=1}^{\infty} \frac{1}{n^2} < \infty , \end{equation*}

so 8.54(a) implies that the unordered sum

\begin{equation*} f = \sum_{k \in \Gamma} \alpha_k e_k \end{equation*}

converges to some \(f \in V\) (here we use that \(V\) is a Hilbert space, hence complete).

Next we compute \(\langle f, e_j \rangle\) for \(j \in \Gamma\). Let \(\varepsilon > 0\). By Definition 8.53 there is a finite set \(\Omega \subseteq \Gamma\) such that \(\| f - \sum_{i \in \Omega^{\prime}} \alpha_i e_i \| < \varepsilon\) for every finite set \(\Omega^{\prime}\) with \(\Omega \subseteq \Omega^{\prime} \subseteq \Gamma\). Take \(\Omega^{\prime} = \Omega \cup \{ j \}\). Then \(\langle \sum_{i \in \Omega^{\prime}} \alpha_i e_i , e_j \rangle = \alpha_j\) by orthonormality, and so by the Cauchy–Schwarz inequality (8.11) and \(\|e_j\| = 1\),

\begin{equation*} \big| \langle f, e_j \rangle - \alpha_j \big| = \Big| \Big\langle f - \sum_{i \in \Omega^{\prime}} \alpha_i e_i , \, e_j \Big\rangle \Big| \le \Big\| f - \sum_{i \in \Omega^{\prime}} \alpha_i e_i \Big\| < \varepsilon . \end{equation*}

Because \(\varepsilon > 0\) was arbitrary, \(\langle f, e_j \rangle = \alpha_j\) for every \(j \in \Gamma\). In particular

\begin{equation*} \langle f, e_{k_n} \rangle = \frac{1}{n} \neq 0 \quad \text{for every } n \in \mathbf{Z}^+ . \end{equation*}

Now we use the assumption that \(\{e_k\}_{k \in \Gamma}\) spans \(V\). Because \(f \in V\), there exist a finite set \(\Omega_0 \subseteq \Gamma\) and a family \(\{\beta_j\}_{j \in \Omega_0}\) in \(\mathbf{F}\) with

\begin{equation*} f = \sum_{j \in \Omega_0} \beta_j e_j . \end{equation*}

Since \(\Omega_0\) is finite and \(k_1, k_2, \ldots\) are distinct elements of \(\Gamma\), there exists \(n \in \mathbf{Z}^+\) with \(k_n \notin \Omega_0\). For that \(n\), orthonormality gives

\begin{equation*} \langle f, e_{k_n} \rangle = \sum_{j \in \Omega_0} \beta_j \langle e_j, e_{k_n} \rangle = 0 , \end{equation*}

because every \(j \in \Omega_0\) differs from \(k_n\). This contradicts \(\langle f, e_{k_n} \rangle = \frac{1}{n} \neq 0\).

This contradiction shows that no orthonormal family in an infinite-dimensional Hilbert space \(V\) can be a basis of \(V\).

Remark: the argument shows more, namely that if \(\{e_k\}_{k \in \Gamma}\) is an orthonormal family in a Hilbert space and \(\Gamma\) is infinite, then \(\operatorname{span}\{e_k\}_{k \in \Gamma}\) is a proper subset of \(\overline{\operatorname{span}\{e_k\}_{k \in \Gamma}}\), since the element \(f\) constructed above lies in the closure of the span by 8.58(a) but not in the span. In particular an orthonormal basis of an infinite-dimensional Hilbert space is never a basis in the sense of 6.54.

Problem (8C.10)

(a) Show that the orthonormal family given in the first bullet point of Example 8.51 is an orthonormal basis of \(\ell^2\).

(b) Show that the orthonormal family given in the second bullet point of Example 8.51 is an orthonormal basis of \(\ell^2(\Gamma)\).

(c) Show that the orthonormal family given in the fourth bullet point of Example 8.51 is not an orthonormal basis of \(L^2\big([0,1)\big)\).

(d) Show that the orthonormal family given in the fifth bullet point of Example 8.51 is not an orthonormal basis of \(L^2(\mathbf{R})\).

Solution

We begin with a criterion that will be used in all four parts.

Lemma. Suppose \(\{e_k\}_{k \in \Gamma}\) is an orthonormal family in a Hilbert space \(V\). Then \(\{e_k\}_{k \in \Gamma}\) is an orthonormal basis of \(V\) if and only if the only \(f \in V\) satisfying \(\langle f, e_k \rangle = 0\) for all \(k \in \Gamma\) is \(f = 0\).

Proof of Lemma. Let \(U = \operatorname{span}\{e_k\}_{k \in \Gamma}\), a subspace of \(V\). Every element of \(U\) is a finite linear combination \(\sum_{j \in \Omega} \alpha_j e_j\) with \(\Omega \subseteq \Gamma\) finite, and by additivity and conjugate homogeneity of the inner product in its second slot,

\begin{equation*} \Big\langle f, \sum_{j \in \Omega} \alpha_j e_j \Big\rangle = \sum_{j \in \Omega} \overline{\alpha_j} \, \langle f, e_j \rangle . \end{equation*}

Hence \(\langle f, u \rangle = 0\) for all \(u \in U\) if and only if \(\langle f, e_k \rangle = 0\) for all \(k \in \Gamma\) (for the reverse implication take \(u = e_k\)). In other words,

\begin{equation*} U^{\perp} = \{ f \in V : \langle f, e_k \rangle = 0 \text{ for all } k \in \Gamma \}. \end{equation*}

By 8.42 we have \(\overline{U} = V\) if and only if \(U^{\perp} = \{0\}\). Since by Definition 8.61 the family is an orthonormal basis exactly when \(\overline{U} = V\), the lemma follows. (Alternatively, the “only if” direction is immediate from Parseval’s identity 8.63(c).)

(a) Here \(V = \ell^2\) and, for \(k \in \mathbf{Z}^+\), \(e_k\) is the sequence whose \(k\)th coordinate is \(1\) and whose other coordinates are \(0\). Suppose \(f = (a_1, a_2, \ldots) \in \ell^2\) satisfies \(\langle f, e_k \rangle = 0\) for every \(k \in \mathbf{Z}^+\). The standard inner product on \(\ell^2\) gives

\begin{equation*} \langle f, e_k \rangle = \sum_{n=1}^{\infty} a_n \overline{e_k(n)} = a_k . \end{equation*}

Thus \(a_k = 0\) for every \(k\), so \(f = 0\). By the Lemma, \(\{e_k\}_{k \in \mathbf{Z}^+}\) is an orthonormal basis of \(\ell^2\).

(b) Here \(\Gamma\) is a nonempty set, \(\ell^2(\Gamma) = L^2(\mu)\) with \(\mu\) counting measure on \(\Gamma\), and for \(k \in \Gamma\) the function \(e_k : \Gamma \to \mathbf{F}\) is defined by \(e_k(j) = 1\) if \(j = k\) and \(e_k(j) = 0\) if \(j \neq k\); in other words \(e_k = \chi_{\{k\}}\).

Suppose \(f \in \ell^2(\Gamma)\) satisfies \(\langle f, e_k \rangle = 0\) for every \(k \in \Gamma\). Fix \(k \in \Gamma\). The function \(f \overline{e_k}\) vanishes on \(\Gamma \setminus \{k\}\) and equals \(f(k)\) at \(k\), so it is the simple function \(f(k) \chi_{\{k\}}\), whose integral with respect to counting measure is \(f(k) \mu(\{k\}) = f(k)\). Hence

\begin{equation*} 0 = \langle f, e_k \rangle = \int_{\Gamma} f \overline{e_k} \, d\mu = f(k). \end{equation*}

Thus \(f(k) = 0\) for every \(k \in \Gamma\), so \(f\) is the zero function, and in particular \(f = 0\) as an element of \(L^2(\mu)\). (With counting measure the only set of measure \(0\) is the empty set, so there is no distinction here between vanishing everywhere and vanishing almost everywhere.) By the Lemma, \(\{e_k\}_{k \in \Gamma}\) is an orthonormal basis of \(\ell^2(\Gamma)\). Part (a) is the special case \(\Gamma = \mathbf{Z}^+\).

(c) Here, for each nonnegative integer \(k\), the function \(e_k : [0,1) \to \mathbf{F}\) is defined by

\begin{equation*} e_k(x) = \begin{cases} 1 & \text{if } x \in \big[ \tfrac{n-1}{2^k}, \tfrac{n}{2^k} \big) \text{ for some odd integer } n, \\ -1 & \text{if } x \in \big[ \tfrac{n-1}{2^k}, \tfrac{n}{2^k} \big) \text{ for some even integer } n. \end{cases} \end{equation*}

Note \(e_0 \equiv 1\) on \([0,1)\), that each \(e_k\) is real valued with \(|e_k| = 1\), and hence \(e_k^2 \equiv 1\).

We first record a computation.

Fact. Suppose \(k > m \ge 0\) are integers and \(j \in \{0, 1, \ldots, 2^m - 1\}\). Then

\begin{equation*} \int_{[\, j2^{-m},\, (j+1)2^{-m})} e_k = 0 . \end{equation*}

Indeed, the interval \([\, j2^{-m}, (j+1)2^{-m})\) is the disjoint union of the \(2^{k-m}\) intervals \(\big[ \tfrac{n-1}{2^k}, \tfrac{n}{2^k} \big)\) for \(n = j2^{k-m} + 1, \ldots, (j+1)2^{k-m}\). These are \(2^{k-m}\) consecutive integers \(n\), and \(2^{k-m}\) is even because \(k > m\); hence exactly half of them are odd and half are even. Each of these intervals has length \(2^{-k}\), and \(e_k\) equals \(1\) on those with \(n\) odd and \(-1\) on those with \(n\) even. So the integral is \(2^{-k}\big( \tfrac{1}{2} 2^{k-m} - \tfrac{1}{2} 2^{k-m}\big) = 0\), proving the Fact.

Now define \(f : [0,1) \to \mathbf{R}\) by \(f = e_1 e_2\). Reading off the definitions, on the four quarter-intervals \([0,\tfrac14)\), \([\tfrac14,\tfrac12)\), \([\tfrac12,\tfrac34)\), \([\tfrac34,1)\) we have

\begin{equation*} e_0 = (1,1,1,1), \qquad e_1 = (1,1,-1,-1), \qquad e_2 = (1,-1,1,-1), \end{equation*}

and therefore

\begin{equation*} f = e_1 e_2 = (1,-1,-1,1), \end{equation*}

meaning \(f\) takes the listed value on the corresponding quarter-interval. Since \(|f| = 1\) on \([0,1)\), we have \(f \in L^2\big([0,1)\big)\) and \(\|f\|_2 = 1\); in particular \(f \neq 0\).

We claim \(\langle f, e_k \rangle = 0\) for every nonnegative integer \(k\). Because everything in sight is real valued, \(\langle f, e_k \rangle = \int_0^1 f e_k\).

For \(k \in \{0, 1, 2\}\) we use the four-quarter description above. Each of \(f e_0\), \(f e_1\), \(f e_2\) is constant on each quarter-interval, with values

\begin{equation*} f e_0 = (1,-1,-1,1), \qquad f e_1 = (1,-1,1,-1), \qquad f e_2 = (1,1,-1,-1), \end{equation*}

and each of these three lists sums to \(0\). Since each quarter-interval has length \(\tfrac14\), each of \(\int_0^1 f e_0\), \(\int_0^1 f e_1\), \(\int_0^1 f e_2\) equals \(0\).

For \(k \ge 3\), note that \(f\) is constant on each quarter-interval, say \(f = c_j\) on \(I_j = [\, j 2^{-2}, (j+1)2^{-2})\) for \(j \in \{0,1,2,3\}\). Applying the Fact with \(m = 2 < k\) gives

\begin{equation*} \int_0^1 f e_k = \sum_{j=0}^{3} c_j \int_{I_j} e_k = 0 . \end{equation*}

So \(f\) is a nonzero element of \(L^2\big([0,1)\big)\) orthogonal to every \(e_k\). By the Lemma, \(\{e_k\}_{k \in \{0,1,\ldots\}}\) is not an orthonormal basis of \(L^2\big([0,1)\big)\). (Equivalently: Parseval’s identity 8.63(c) would force \(1 = \|f\|^2 = \sum_k |\langle f, e_k\rangle|^2 = 0\).)

(d) Here, for \(k\) a nonnegative integer and \(m \in \mathbf{Z}\), the function \(e_{k,m} : \mathbf{R} \to \mathbf{F}\) is the translate by \(m\) of the function \(e_k\) of part (c), extended by \(0\); that is,

\begin{equation*} e_{k,m}(x) = \begin{cases} e_k(x - m) & \text{if } x \in [m, m+1), \\ 0 & \text{if } x \notin [m, m+1). \end{cases} \end{equation*}

(This is exactly the fifth bullet point of Example 8.51: the condition \(x \in \big[ m + \tfrac{n-1}{2^k}, m + \tfrac{n}{2^k}\big)\) with \(n \in [1, 2^k]\) an integer says precisely that \(x - m \in \big[\tfrac{n-1}{2^k}, \tfrac{n}{2^k}\big) \subseteq [0,1)\).)

Let \(f : \mathbf{R} \to \mathbf{R}\) be the extension by \(0\) of the function \(f\) from part (c):

\begin{equation*} f = \chi_{[0, 1/4)} - \chi_{[1/4, 1/2)} - \chi_{[1/2, 3/4)} + \chi_{[3/4, 1)} . \end{equation*}

Then \(|f| = \chi_{[0,1)}\), so \(f \in L^2(\mathbf{R})\) with \(\|f\|_2 = 1\); in particular \(f \neq 0\).

Fix a nonnegative integer \(k\) and \(m \in \mathbf{Z}\).

If \(m \neq 0\), then \([m, m+1) \cap [0,1) = \varnothing\) because \(m\) is an integer. Since \(f\) vanishes outside \([0,1)\) and \(e_{k,m}\) vanishes outside \([m,m+1)\), the product \(f \overline{e_{k,m}}\) is identically \(0\), so \(\langle f, e_{k,m} \rangle = 0\).

If \(m = 0\), then \(e_{k,0}\) agrees with \(e_k\) on \([0,1)\) and vanishes elsewhere, while \(f\) vanishes outside \([0,1)\). Hence

\begin{equation*} \langle f, e_{k,0} \rangle = \int_{\mathbf{R}} f \, e_{k,0} = \int_0^1 f e_k = 0 \end{equation*}

by part (c).

Thus \(f\) is a nonzero element of \(L^2(\mathbf{R})\) orthogonal to every element of the family \(\{e_{k,m}\}_{(k,m) \in \{0,1,\ldots\} \times \mathbf{Z}}\). By the Lemma, this family is not an orthonormal basis of \(L^2(\mathbf{R})\).

Problem (8C.11)

Suppose \(\mu\) is a \(\sigma\)-finite measure on \((X, \mathcal{S})\) and \(\nu\) is a \(\sigma\)-finite measure on \((Y, \mathcal{T})\). Suppose also that \(\{e_j\}_{j \in \Omega}\) is an orthonormal basis of \(L^2(\mu)\) and \(\{f_k\}_{k \in \Gamma}\) is an orthonormal basis of \(L^2(\nu)\) for some countable set \(\Gamma\). For \(j \in \Omega\) and \(k \in \Gamma\), define \(g_{j,k} : X \times Y \to \mathbf{F}\) by

\begin{equation*} g_{j,k}(x,y) = e_j(x) f_k(y). \end{equation*}

Prove that \(\{g_{j,k}\}_{j \in \Omega, \, k \in \Gamma}\) is an orthonormal basis of \(L^2(\mu \times \nu)\).

Solution

Throughout, the hypothesis that \(\mu\) and \(\nu\) are \(\sigma\)-finite is exactly what licenses every application below of Tonelli’s Theorem (5.28) and Fubini’s Theorem (5.32), and it also makes \(\mu \times \nu\) a well-defined measure on \(\mathcal{S} \otimes \mathcal{T}\); without it none of what follows is available. Note also that \(L^2(\mu \times \nu)\) is a Hilbert space (first bullet point of 8.22), so the results of Section 8B may be applied to it. Finally, 5.28 and 5.32 are stated in MIRA for functions with values in \([0, \infty]\) and \([-\infty, \infty]\) respectively; when \(\mathbf{F} = \mathbf{C}\) we apply them to the real and imaginary parts (and, for Tonelli, to the positive and negative parts of those), as usual.

We use the criterion proved as the Lemma in Exercise 10: an orthonormal family \(\{u_i\}_{i \in I}\) in a Hilbert space \(V\) is an orthonormal basis of \(V\) if and only if the only \(h \in V\) with \(\langle h, u_i \rangle = 0\) for all \(i \in I\) is \(h = 0\). (Proof: if \(U = \operatorname{span}\{u_i\}_{i \in I}\), then \(U^{\perp}\) is exactly the set of such \(h\), and by 8.42 we have \(\overline{U} = V\) if and only if \(U^{\perp} = \{0\}\).)

Step 1: each \(g_{j,k}\) lies in \(L^2(\mu \times \nu)\) and has norm \(1\).

The function \((x,y) \mapsto e_j(x)\) is \(\mathcal{S} \otimes \mathcal{T}\)-measurable, because for each Borel set \(B \subseteq \mathbf{F}\) its inverse image is \(e_j^{-1}(B) \times Y\), which lies in \(\mathcal{S} \otimes \mathcal{T}\); similarly \((x,y) \mapsto f_k(y)\) is \(\mathcal{S} \otimes \mathcal{T}\)-measurable. A product of measurable functions is measurable, so \(g_{j,k}\) is \(\mathcal{S} \otimes \mathcal{T}\)-measurable. Now \(|g_{j,k}(x,y)|^2 = |e_j(x)|^2 |f_k(y)|^2 \ge 0\), so Tonelli’s Theorem (5.28) applies and gives

\begin{equation*} \int_{X \times Y} |g_{j,k}|^2 \, d(\mu \times \nu) = \int_X \Big( \int_Y |e_j(x)|^2 |f_k(y)|^2 \, d\nu(y) \Big) d\mu(x) = \|e_j\|^2 \, \|f_k\|^2 = 1 . \end{equation*}

Thus \(g_{j,k} \in L^2(\mu \times \nu)\) and \(\|g_{j,k}\| = 1\).

Step 2: the family is orthonormal.

Let \(j, j^{\prime} \in \Omega\) and \(k, k^{\prime} \in \Gamma\). Since \(g_{j,k}\) and \(g_{j^{\prime},k^{\prime}}\) both lie in \(L^2(\mu \times \nu)\), the Cauchy–Schwarz inequality (8.11), or equivalently Holder’s inequality (7.9), shows \(g_{j,k}\overline{g_{j^{\prime},k^{\prime}}} \in L^1(\mu \times \nu)\). Hence Fubini’s Theorem (5.32) applies and yields

\begin{equation*} \begin{aligned} \langle g_{j,k}, g_{j^{\prime},k^{\prime}} \rangle &= \int_{X \times Y} e_j(x) f_k(y) \overline{e_{j^{\prime}}(x)} \, \overline{f_{k^{\prime}}(y)} \, d(\mu \times \nu)(x,y) \\ &= \int_X e_j(x) \overline{e_{j^{\prime}}(x)} \Big( \int_Y f_k(y) \overline{f_{k^{\prime}}(y)} \, d\nu(y) \Big) d\mu(x) \\ &= \langle e_j, e_{j^{\prime}} \rangle \, \langle f_k, f_{k^{\prime}} \rangle . \end{aligned} \end{equation*}

If \((j,k) = (j^{\prime},k^{\prime})\) this equals \(1\); if \((j,k) \neq (j^{\prime},k^{\prime})\) then either \(j \neq j^{\prime}\), making \(\langle e_j, e_{j^{\prime}}\rangle = 0\), or \(k \neq k^{\prime}\), making \(\langle f_k, f_{k^{\prime}} \rangle = 0\). Either way the inner product is \(0\). Thus \(\{g_{j,k}\}_{(j,k) \in \Omega \times \Gamma}\) is an orthonormal family in \(L^2(\mu \times \nu)\).

Step 3: nothing nonzero is orthogonal to the whole family.

Suppose \(h \in L^2(\mu \times \nu)\) satisfies \(\langle h, g_{j,k}\rangle = 0\) for all \(j \in \Omega\) and \(k \in \Gamma\). Fix an \(\mathcal{S} \otimes \mathcal{T}\)-measurable representative of \(h\).

Because \(|h|^2\) is a nonnegative \(\mathcal{S} \otimes \mathcal{T}\)-measurable function, Tonelli’s Theorem (5.28) says that

\begin{equation*} \varphi(x) = \int_Y |h(x,y)|^2 \, d\nu(y) \end{equation*}

defines an \(\mathcal{S}\)-measurable function \(\varphi : X \to [0, \infty]\) with

\begin{equation*} \int_X \varphi \, d\mu = \int_{X \times Y} |h|^2 \, d(\mu \times \nu) = \|h\|^2 < \infty . \end{equation*}

Hence \(\varphi(x) < \infty\) for \(\mu\)-almost every \(x\); let \(E \in \mathcal{S}\) be such that \(\mu(X \setminus E) = 0\) and \(\varphi(x) < \infty\) for all \(x \in E\). For \(x \in X\) write \(h_x\) for the cross section \(y \mapsto h(x,y)\), which is \(\mathcal{T}\)-measurable by 5.9. For \(x \in E\) we thus have \(h_x \in L^2(\nu)\) with \(\|h_x\|^2 = \varphi(x)\).

For \(k \in \Gamma\) define \(H_k : X \to \mathbf{F}\) by

\begin{equation*} H_k(x) = \langle h_x, f_k \rangle = \int_Y h(x,y) \overline{f_k(y)} \, d\nu(y) \quad \text{for } x \in E, \qquad H_k(x) = 0 \quad \text{for } x \in X \setminus E . \end{equation*}

This is well defined: by the Cauchy–Schwarz inequality (8.11) in \(L^2(\nu)\),

\begin{equation*} \int_Y |h(x,y) \overline{f_k(y)}| \, d\nu(y) \le \|h_x\| \, \|f_k\| = \varphi(x)^{1/2} < \infty \end{equation*}

for \(x \in E\), and the same estimate gives

\begin{equation*} |H_k(x)| \le \varphi(x)^{1/2} \qquad \text{for all } x \in X \end{equation*}

(trivially at points of \(X \setminus E\), where both sides may be read as \(0 \le \varphi(x)^{1/2}\)). Moreover \(H_k\) is \(\mathcal{S}\)-measurable: the function \((x,y) \mapsto h(x,y)\overline{f_k(y)}\) is \(\mathcal{S} \otimes \mathcal{T}\)-measurable, and applying 5.28 to the positive and negative parts of its real and imaginary parts produces four \(\mathcal{S}\)-measurable functions \(X \to [0, \infty]\), each finite on \(E\) by the displayed bound; on \(E\) the function \(H_k\) is the corresponding combination of them, and \(E \in \mathcal{S}\), so multiplying that combination by \(\chi_E\) gives an \(\mathcal{S}\)-measurable function on \(X\) equal to \(H_k\). Since \(|H_k|^2 \le \varphi\) and \(\varphi \in L^1(\mu)\), we get \(H_k \in L^2(\mu)\).

Now we use each of the two orthonormal bases in turn.

First, since \(\{f_k\}_{k \in \Gamma}\) is an orthonormal basis of \(L^2(\nu)\), Parseval’s identity 8.63(c) applied to \(h_x \in L^2(\nu)\) gives, for every \(x \in E\),

\begin{equation*} \varphi(x) = \|h_x\|^2 = \sum_{k \in \Gamma} |\langle h_x, f_k \rangle|^2 = \sum_{k \in \Gamma} |H_k(x)|^2 . \end{equation*}

Second, fix \(k \in \Gamma\) and \(j \in \Omega\). As in Step 2, \(h \overline{g_{j,k}} \in L^1(\mu \times \nu)\) because \(h\) and \(g_{j,k}\) lie in \(L^2(\mu \times \nu)\). So Fubini’s Theorem (5.32) gives

\begin{equation*} 0 = \langle h, g_{j,k} \rangle = \int_X \Big( \int_Y h(x,y) \overline{f_k(y)} \, d\nu(y) \Big) \overline{e_j(x)} \, d\mu(x) = \int_X H_k \overline{e_j} \, d\mu = \langle H_k, e_j \rangle , \end{equation*}

where the inner integral is \(H_k(x)\) for \(\mu\)-almost every \(x\) (namely for \(x \in E\)), which is all that is needed for the integral over \(X\). Thus \(\langle H_k, e_j \rangle = 0\) for every \(j \in \Omega\). Because \(\{e_j\}_{j \in \Omega}\) is an orthonormal basis of \(L^2(\mu)\), Parseval’s identity 8.63(c) now gives

\begin{equation*} \|H_k\|^2 = \sum_{j \in \Omega} |\langle H_k, e_j \rangle|^2 = 0 , \end{equation*}

so \(H_k = 0\) \(\mu\)-almost everywhere on \(X\).

Here is where the countability of \(\Gamma\) enters. For each \(k \in \Gamma\) the set \(N_k = \{x \in X : H_k(x) \neq 0\}\) belongs to \(\mathcal{S}\) (because \(H_k\) is \(\mathcal{S}\)-measurable) and satisfies \(\mu(N_k) = 0\). Let \(N = \bigcup_{k \in \Gamma} N_k\). Since \(\Gamma\) is countable, \(N \in \mathcal{S}\) and \(\mu(N) = 0\) by countable subadditivity. For \(x \in E \setminus N\) we have \(H_k(x) = 0\) for every \(k \in \Gamma\), so the first display above gives

\begin{equation*} \varphi(x) = \sum_{k \in \Gamma} |H_k(x)|^2 = 0 . \end{equation*}

Thus \(\varphi = 0\) \(\mu\)-almost everywhere on \(X\) (recall \(\mu(X \setminus E) = 0\)), and therefore

\begin{equation*} \|h\|^2 = \int_{X \times Y} |h|^2 \, d(\mu \times \nu) = \int_X \varphi \, d\mu = 0 . \end{equation*}

Hence \(h = 0\) in \(L^2(\mu \times \nu)\).

By Steps 1–3 and the criterion recalled at the start, \(\{g_{j,k}\}_{j \in \Omega, \, k \in \Gamma}\) is an orthonormal basis of \(L^2(\mu \times \nu)\).

Problem (8C.12)

Prove the converse of Parseval’s identity. More specifically, prove that if \(\{e_k\}_{k \in \Gamma}\) is an orthonormal family in a Hilbert space \(V\) and

\begin{equation*} \|f\|^2 = \sum_{k \in \Gamma} |\langle f, e_k \rangle|^2 \end{equation*}

for every \(f \in V\), then \(\{e_k\}_{k \in \Gamma}\) is an orthonormal basis of \(V\).

Solution

Let

\begin{equation*} U = \operatorname{span}\{e_k\}_{k \in \Gamma}, \end{equation*}

which is a subspace of \(V\). By Definition 8.61, we must show \(\overline{U} = V\).

First we identify \(U^{\perp}\). Suppose \(f \in V\). Every element of \(U\) has the form \(\sum_{j \in \Omega} \alpha_j e_j\) for some finite \(\Omega \subseteq \Gamma\) and some family \(\{\alpha_j\}_{j \in \Omega}\) in \(\mathbf{F}\), and

\begin{equation*} \Big\langle f, \sum_{j \in \Omega} \alpha_j e_j \Big\rangle = \sum_{j \in \Omega} \overline{\alpha_j} \, \langle f, e_j \rangle \end{equation*}

by the conjugate homogeneity and additivity of the inner product in its second slot. Hence if \(\langle f, e_k \rangle = 0\) for every \(k \in \Gamma\), then \(\langle f, u \rangle = 0\) for every \(u \in U\); and conversely, taking \(u = e_k\) shows the other implication. Therefore

\begin{equation*} U^{\perp} = \{ f \in V : \langle f, e_k \rangle = 0 \text{ for all } k \in \Gamma \}. \end{equation*}

Now suppose \(f \in U^{\perp}\). Then \(\langle f, e_k \rangle = 0\) for every \(k \in \Gamma\), so the hypothesis gives

\begin{equation*} \|f\|^2 = \sum_{k \in \Gamma} |\langle f, e_k \rangle|^2 = 0 , \end{equation*}

which forces \(f = 0\) by the definition of a norm. Hence

\begin{equation*} U^{\perp} = \{0\}. \end{equation*}

Because \(V\) is a Hilbert space and \(U\) is a subspace of \(V\), the criterion 8.42 (a subspace of a Hilbert space is dense if and only if its orthogonal complement is \(\{0\}\)) now gives

\begin{equation*} \overline{\operatorname{span}\{e_k\}_{k \in \Gamma}} = \overline{U} = V . \end{equation*}

Since \(\{e_k\}_{k \in \Gamma}\) is by hypothesis an orthonormal family, Definition 8.61 says precisely that \(\{e_k\}_{k \in \Gamma}\) is an orthonormal basis of \(V\), as desired.

Two remarks. First, the hypothesis is needed only for \(f\) in \(U^{\perp}\), so the following weaker statement suffices: if \(\|f\|^2 = \sum_{k \in \Gamma}|\langle f, e_k\rangle|^2\) whenever \(f\) is orthogonal to every \(e_k\), then \(\{e_k\}_{k \in \Gamma}\) is an orthonormal basis. Second, completeness of \(V\) is essential, since 8.42 depends on the orthogonal decomposition theory of 8.40 and 8.41, which uses the existence of the orthogonal projection \(P_U\) onto a closed subspace (8.37), a result that requires completeness.

Problem (8C.13)

(a) Show that the Hilbert space \(L^2([0,1])\) is separable.

(b) Show that the Hilbert space \(L^2(\mathbf{R})\) is separable.

(c) Show that the Banach space \(\ell^{\infty}\) is not separable.

Solution

Let \(\mathbf{Q}_{\mathbf{F}}\) denote \(\mathbf{Q}\) if \(\mathbf{F} = \mathbf{R}\) and \(\{ r + is : r, s \in \mathbf{Q}\}\) if \(\mathbf{F} = \mathbf{C}\); in either case \(\mathbf{Q}_{\mathbf{F}}\) is a countable dense subset of \(\mathbf{F}\).

For a Lebesgue measurable set \(J \subseteq \mathbf{R}\) let

\begin{equation*} D_J = \Big\{ \sum_{j=1}^{n} q_j \, \chi_{(c_j, d_j)} \big|_J \; : \; n \in \mathbf{Z}^+ , \; q_j \in \mathbf{Q}_{\mathbf{F}}, \; c_j, d_j \in \mathbf{Q}, \; c_j < d_j \Big\} . \end{equation*}

Each element of \(D_J\) is a bounded measurable function vanishing outside a bounded set, hence lies in \(L^2(J)\). The set \(D_J\) is countable, because the map sending a finite sequence \(\big((q_1,c_1,d_1), \ldots, (q_n,c_n,d_n)\big)\) in the countable set \(\mathbf{Q}_{\mathbf{F}} \times \mathbf{Q} \times \mathbf{Q}\) to the corresponding function is a surjection onto \(D_J\), and the set of finite sequences from a countable set is countable.

We prove three preliminary claims. In all of them \(J\) is a Lebesgue measurable subset of \(\mathbf{R}\).

Claim 1 (simple functions are dense). If \(f \in L^2(J)\) and \(\varepsilon > 0\), then there is a simple Lebesgue measurable function \(s\) on \(J\) with \(s \in L^2(J)\) and \(\|f - s\|_2 < \varepsilon\).

Proof. Write \(f = u + iv\) with \(u, v\) real valued and measurable (if \(\mathbf{F} = \mathbf{R}\) take \(v = 0\)); then \(|u| \le |f|\) and \(|v| \le |f|\), so \(u, v \in L^2(J)\). By 2.89 there is a sequence \(u_1, u_2, \ldots\) of simple measurable real valued functions on \(J\) with \(|u_k| \le |u|\) and \(u_k \to u\) pointwise. Then \(|u - u_k|^2 \le 4 u^2\), and \(4u^2 \in L^1(J)\) because \(u \in L^2(J)\), and \(|u - u_k|^2 \to 0\) pointwise. The Dominated Convergence Theorem (3.31) gives \(\int_J |u - u_k|^2 \to 0\), that is, \(\|u - u_k\|_2 \to 0\). Similarly there are simple \(v_k\) with \(\|v - v_k\|_2 \to 0\). Taking \(s = u_k + i v_k\) for large \(k\) gives \(\|f - s\|_2 \le \|u - u_k\|_2 + \|v - v_k\|_2 < \varepsilon\), and \(s\) is simple and measurable. \(\square\)

Claim 2 (measurable sets of finite measure are almost finite unions of rational intervals). If \(A \subseteq \mathbf{R}\) is Lebesgue measurable with \(|A| < \infty\) and \(\delta > 0\), then there is a set \(W\) that is a finite union of pairwise disjoint bounded open intervals with rational endpoints such that \(|A \triangle W| < \delta\), where \(A \triangle W = (A \setminus W) \cup (W \setminus A)\).

Proof. By part (e) of 2.71 there is an open set \(G \supseteq A\) with \(|G \setminus A| < \delta/4\); in particular \(|G| < \infty\). Every open subset of \(\mathbf{R}\) is a countable union of pairwise disjoint nonempty open intervals, say \(G = \bigcup_{j=1}^{\infty} I_j\) (a finite union, or the empty union when \(G = \varnothing\), is handled the same way, and then one takes \(W\) to be a union of fewer intervals or \(W = \varnothing\)). Each \(I_j\) is bounded since \(|G| < \infty\), and \(\sum_j |I_j| = |G| < \infty\) by countable additivity, so there is \(N \in \mathbf{Z}^+\) with \(\sum_{j > N} |I_j| < \delta/4\); put \(U = \bigcup_{j \le N} I_j\), so \(|G \setminus U| = \sum_{j > N} |I_j| < \delta/4\). Writing \(I_j = (\alpha_j, \beta_j)\), which is nonempty and bounded and hence has \(-\infty < \alpha_j < \beta_j < \infty\), choose rationals \(c_j, d_j\) with \(\alpha_j < c_j < d_j < \beta_j\) and \(|I_j| - (d_j - c_j) < \delta/(4N)\), and set \(W = \bigcup_{j \le N} (c_j, d_j)\). The intervals \((c_j, d_j)\) for \(j \le N\) are pairwise disjoint because \((c_j, d_j) \subseteq I_j\) and the \(I_j\) are pairwise disjoint. Also \(W \subseteq U\) and \(|U \setminus W| = \sum_{j \le N} \big( |I_j| - (d_j - c_j) \big) < \delta/4\). Consequently

\begin{equation*} |A \setminus W| \le |A \setminus G| + |G \setminus U| + |U \setminus W| < 0 + \tfrac{\delta}{4} + \tfrac{\delta}{4}, \end{equation*}

and \(|W \setminus A| \le |G \setminus A| < \delta/4\). Hence \(|A \triangle W| < \delta\). \(\square\)

Claim 3. Suppose \(J \subseteq \mathbf{R}\) is measurable with \(|J| < \infty\). Then \(D_J\) is dense in \(L^2(J)\).

Proof. Let \(f \in L^2(J)\) and \(\varepsilon > 0\). By Claim 1 there is a simple measurable \(s \in L^2(J)\) with \(\|f - s\|_2 < \varepsilon/3\). Write \(s = \sum_{k=1}^{n} a_k \chi_{A_k}\) where \(a_1, \ldots, a_n\) are nonzero elements of \(\mathbf{F}\) and \(A_1, \ldots, A_n\) are disjoint measurable subsets of \(J\); note \(|A_k| \le |J| < \infty\).

Choose \(q_1, \ldots, q_n \in \mathbf{Q}_{\mathbf{F}}\) with \(|a_k - q_k| < \varepsilon / \big(3n(1 + |J|)^{1/2}\big)\) for each \(k\). Then, using \(\|\chi_{A_k}\|_2 = |A_k|^{1/2} \le |J|^{1/2}\) and the triangle inequality,

\begin{equation*} \Big\| s - \sum_{k=1}^{n} q_k \chi_{A_k} \Big\|_2 \le \sum_{k=1}^{n} |a_k - q_k| \, |A_k|^{1/2} < \frac{\varepsilon}{3} . \end{equation*}

Now let \(\delta > 0\) satisfy \(\delta^{1/2} \sum_{k=1}^{n} |q_k| < \varepsilon/3\). By Claim 2, for each \(k\) there is a finite union \(W_k\) of bounded open intervals with rational endpoints such that \(|A_k \triangle W_k| < \delta\). Since \(|\chi_{A_k} - \chi_{W_k \cap J}| \le \chi_{A_k \triangle W_k}\) on \(J\) (recall \(A_k \subseteq J\)),

\begin{equation*} \big\| \chi_{A_k} - \chi_{W_k}|_J \big\|_2 \le |A_k \triangle W_k|^{1/2} < \delta^{1/2}, \end{equation*}

and therefore

\begin{equation*} \Big\| \sum_{k=1}^{n} q_k \chi_{A_k} - \sum_{k=1}^{n} q_k \chi_{W_k}\big|_J \Big\|_2 \le \sum_{k=1}^{n} |q_k| \, \delta^{1/2} < \frac{\varepsilon}{3} . \end{equation*}

Finally, each \(W_k\) is a finite union of pairwise disjoint bounded open intervals with rational endpoints, so \(\chi_{W_k}\) is the sum of the characteristic functions of those intervals (disjointness is what makes this an equality rather than an inequality). Hence \(\sum_{k} q_k \chi_{W_k}|_J\) is a finite sum of terms \(q \chi_{(c,d)}|_J\) with \(q \in \mathbf{Q}_{\mathbf{F}}\) and \(c < d\) rational; that is, it belongs to \(D_J\). (Degenerate cases cause no trouble: \(0 \in D_J\), since \(0 \in \mathbf{Q}_{\mathbf{F}}\) and one may take the single term \(0 \cdot \chi_{(0,1)}\).) Adding the three estimates gives an element of \(D_J\) within \(\varepsilon\) of \(f\). \(\square\)

(a) Take \(J = [0,1]\), so \(|J| = 1 < \infty\). By Claim 3, \(D_{[0,1]}\) is a countable subset of \(L^2([0,1])\) whose closure is all of \(L^2([0,1])\). Hence \(L^2([0,1])\) is separable, by Definition 8.64.

(b) We show \(D_{\mathbf{R}}\), which is countable, is dense in \(L^2(\mathbf{R})\). Let \(f \in L^2(\mathbf{R})\) and \(\varepsilon > 0\). For \(m \in \mathbf{Z}^+\) set \(f_m = f \chi_{[-m,m]}\). Then \(|f - f_m|^2 = |f|^2 \chi_{\mathbf{R} \setminus [-m,m]} \to 0\) pointwise and \(|f - f_m|^2 \le |f|^2 \in L^1(\mathbf{R})\), so the Dominated Convergence Theorem (3.31) gives \(\|f - f_m\|_2 \to 0\). Choose \(m\) with \(\|f - f_m\|_2 < \varepsilon/2\).

The restriction of \(f_m\) to \([-m,m]\) lies in \(L^2([-m,m])\), and \(|[-m,m]| = 2m < \infty\), so Claim 3 provides \(g \in D_{[-m,m]}\) with \(\|f_m|_{[-m,m]} - g\|_{L^2([-m,m])} < \varepsilon/2\). Let \(\tilde{g} \in D_{\mathbf{R}}\) be the corresponding function on \(\mathbf{R}\), that is, the same finite sum \(\sum_j q_j \chi_{(c_j,d_j)}\) but regarded as a function on \(\mathbf{R}\); replacing each interval \((c_j,d_j)\) by \((c_j,d_j) \cap (-m,m)\) (again an interval with rational endpoints, discarded if empty) we may assume \(\tilde{g}\) vanishes outside \([-m,m]\). Since \(f_m\) also vanishes outside \([-m,m]\),

\begin{equation*} \|f_m - \tilde{g}\|_{L^2(\mathbf{R})} = \big\| f_m|_{[-m,m]} - g \big\|_{L^2([-m,m])} < \frac{\varepsilon}{2}, \end{equation*}

and hence \(\|f - \tilde{g}\|_{L^2(\mathbf{R})} < \varepsilon\). Thus \(D_{\mathbf{R}}\) is dense in \(L^2(\mathbf{R})\), so \(L^2(\mathbf{R})\) is separable.

(c) Recall that \(\ell^{\infty}\) is the Banach space of bounded sequences in \(\mathbf{F}\) with \(\|(a_1, a_2, \ldots)\|_{\infty} = \sup_k |a_k|\).

For each subset \(S \subseteq \mathbf{Z}^+\) let \(x_S \in \ell^{\infty}\) be the sequence whose \(k\)th term is \(1\) if \(k \in S\) and \(0\) if \(k \notin S\). If \(S\) and \(S^{\prime}\) are distinct subsets of \(\mathbf{Z}^+\), then there is \(k\) lying in exactly one of them, so the \(k\)th term of \(x_S - x_{S^{\prime}}\) has absolute value \(1\); since all terms of \(x_S - x_{S^{\prime}}\) have absolute value at most \(1\), we get

\begin{equation*} \|x_S - x_{S^{\prime}}\|_{\infty} = 1 . \end{equation*}

Consider the open balls \(B(x_S, \tfrac12)\) for \(S \subseteq \mathbf{Z}^+\). These are pairwise disjoint: if \(y \in B(x_S, \tfrac12) \cap B(x_{S^{\prime}}, \tfrac12)\) with \(S \ne S^{\prime}\), then

\begin{equation*} 1 = \|x_S - x_{S^{\prime}}\|_{\infty} \le \|x_S - y\|_{\infty} + \|y - x_{S^{\prime}}\|_{\infty} < \tfrac12 + \tfrac12 = 1, \end{equation*}

a contradiction. By Cantor’s theorem the collection of all subsets of \(\mathbf{Z}^+\) is uncountable, so \(\{B(x_S, \tfrac12) : S \subseteq \mathbf{Z}^+\}\) is an uncountable collection of pairwise disjoint nonempty open subsets of \(\ell^{\infty}\).

Now suppose \(C\) were a countable subset of \(\ell^{\infty}\) with \(\overline{C} = \ell^{\infty}\). A dense subset of a metric space intersects every nonempty open set (see the elementary facts following the definition of dense in 6.75), so for each \(S \subseteq \mathbf{Z}^+\) we could choose \(y_S \in C \cap B(x_S, \tfrac12)\). The balls are pairwise disjoint, so \(S \mapsto y_S\) is injective, exhibiting an injection of an uncountable set into the countable set \(C\), a contradiction.

Therefore \(\ell^{\infty}\) has no countable dense subset; that is, \(\ell^{\infty}\) is not separable. (This is the same argument used in 8.66 for \(\ell^2(\Gamma)\) with \(\Gamma\) uncountable and for \(L^{\infty}([0,1])\).)

Problem (8C.14)

Prove that every subspace of a separable normed vector space is separable.

Solution

Suppose \(V\) is a separable normed vector space and \(U\) is a subspace of \(V\). Give \(U\) the norm inherited from \(V\); then \(U\) is a normed vector space, and the metric on \(U\) induced by this norm is the restriction of the metric of \(V\). So we must produce a countable subset of \(U\) whose closure (in \(U\)) is \(U\).

By Definition 8.64 there is a countable subset of \(V\) whose closure equals \(V\). Write this set as \(\{f_n\}_{n \in I}\) where \(I \subseteq \mathbf{Z}^+\) (a countable set can be indexed by a subset of \(\mathbf{Z}^+\)). Since \(V \ne \varnothing\), the set is nonempty, so \(I \ne \varnothing\).

For \(n \in I\) and \(k \in \mathbf{Z}^+\), let

\begin{equation*} B\big(f_n, \tfrac{1}{k}\big) = \Big\{ g \in V : \|g - f_n\| < \tfrac{1}{k} \Big\} \end{equation*}

be the open ball in \(V\) of radius \(1/k\) centered at \(f_n\). Let

\begin{equation*} P = \Big\{ (n,k) \in I \times \mathbf{Z}^+ : B\big(f_n, \tfrac{1}{k}\big) \cap U \ne \varnothing \Big\} , \end{equation*}

and for each \((n,k) \in P\) choose an element

\begin{equation*} u_{n,k} \in B\big(f_n, \tfrac{1}{k}\big) \cap U . \end{equation*}

Set

\begin{equation*} C = \{ u_{n,k} : (n,k) \in P \} . \end{equation*}

Then \(C \subseteq U\), and \(C\) is countable because \(P \subseteq I \times \mathbf{Z}^+ \subseteq \mathbf{Z}^+ \times \mathbf{Z}^+\) is countable and \(C\) is the image of \(P\) under \((n,k) \mapsto u_{n,k}\).

We claim that the closure of \(C\) in \(U\) equals \(U\). Let \(u \in U\) and let \(\varepsilon > 0\). Choose \(k \in \mathbf{Z}^+\) with \(\frac{2}{k} < \varepsilon\). Because \(\{f_n\}_{n \in I}\) is dense in \(V\) and \(u \in V\), there exists \(n \in I\) with

\begin{equation*} \|u - f_n\| < \tfrac{1}{k} . \end{equation*}

Then \(u \in B(f_n, \tfrac{1}{k}) \cap U\), so this intersection is nonempty, that is, \((n,k) \in P\) and \(u_{n,k}\) is defined; moreover \(\|u_{n,k} - f_n\| < \tfrac{1}{k}\) by construction. The triangle inequality now gives

\begin{equation*} \|u - u_{n,k}\| \le \|u - f_n\| + \|f_n - u_{n,k}\| < \tfrac{1}{k} + \tfrac{1}{k} = \tfrac{2}{k} < \varepsilon . \end{equation*}

Thus for every \(u \in U\) and every \(\varepsilon > 0\) there is an element of \(C\) within distance \(\varepsilon\) of \(u\). By the characterization of closure in a metric space (6.7), this says \(u\) is in the closure of \(C\). Hence \(U \subseteq \overline{C}\), and since \(C \subseteq U\) and \(U\) is being regarded as a metric space in its own right, the closure of \(C\) in \(U\) is all of \(U\).

Therefore \(U\) has a countable subset whose closure equals \(U\), so \(U\) is separable.

Two remarks. First, the argument used no algebraic structure at all: exactly the same proof shows that every subset of a separable metric space is separable (with the inherited metric). Second, the construction of the family \(\{u_{n,k}\}_{(n,k) \in P}\) uses the axiom of countable choice, which is standard in this context.

Problem (8C.15)

Suppose \(V\) is an infinite-dimensional Hilbert space. Prove that there does not exist a translation invariant measure on the Borel subsets of \(V\) that assigns positive but finite measure to each open ball in \(V\).

[A subset of \(V\) is called a Borel set if it is in the smallest \(\sigma\)-algebra containing all the open subsets of \(V\). A measure \(\mu\) on the Borel subsets of \(V\) is called translation invariant if \(\mu(f + E) = \mu(E)\) for every \(f \in V\) and every Borel set \(E\) of \(V\).]

Solution

For \(f \in V\) and \(r > 0\) write \(B(f, r) = \{g \in V : \|g - f\| < r\}\) for the open ball of radius \(r\) centered at \(f\). Every such ball is an open subset of \(V\) and hence is a Borel set.

Suppose, seeking a contradiction, that \(\mu\) is a translation invariant measure on the Borel subsets of \(V\) such that \(0 < \mu(B) < \infty\) for every open ball \(B\) in \(V\).

Step 1: \(V\) contains an orthonormal sequence. By 8.75, \(V\) has an orthonormal basis \(\{e_k\}_{k \in \Gamma}\). If \(\Gamma\) were a finite set, then Parseval’s identity [8.63(a)] would give

\begin{equation*} f = \sum_{k \in \Gamma} \langle f, e_k \rangle e_k \in \operatorname{span}\{e_k\}_{k \in \Gamma} \end{equation*}

for every \(f \in V\), making \(V\) a finite-dimensional vector space and contradicting the hypothesis that \(V\) is infinite-dimensional. Hence \(\Gamma\) is infinite. Choose distinct indices \(k_1, k_2, \ldots \in \Gamma\) and set \(u_n = e_{k_n}\). Then \(u_1, u_2, \ldots\) is an orthonormal sequence in \(V\).

Step 2: the \(u_n\) are uniformly separated. If \(m \neq n\), then \(\langle u_m, u_n \rangle = 0\), so by the Pythagorean theorem

\begin{equation*} \|u_m - u_n\|^2 = \|u_m\|^2 + \|u_n\|^2 = 1 + 1 = 2 , \end{equation*}

and thus \(\|u_m - u_n\| = \sqrt{2}\).

Step 3: disjoint balls of a fixed radius. Let \(r = \frac{\sqrt{2}}{2}\). The balls \(B(u_1, r), B(u_2, r), \ldots\) are pairwise disjoint: if \(m \neq n\) and \(h \in B(u_m, r) \cap B(u_n, r)\), then the triangle inequality gives

\begin{equation*} \sqrt{2} = \|u_m - u_n\| \le \|u_m - h\| + \|h - u_n\| < r + r = \sqrt{2}, \end{equation*}

which is a contradiction.

Step 4: all these balls sit inside one ball. If \(h \in B(u_n, r)\), then \(\|h\| \le \|h - u_n\| + \|u_n\| < r + 1\). Hence

\begin{equation*} \bigcup_{n=1}^{\infty} B(u_n, r) \subseteq B(0, 1 + r). \end{equation*}

Step 5: the contradiction. Because \(B(u_n, r) = u_n + B(0, r)\), translation invariance of \(\mu\) gives

\begin{equation*} \mu\bigl(B(u_n, r)\bigr) = \mu\bigl(B(0, r)\bigr) = c \end{equation*}

for every \(n \in \mathbf{Z}^+\), where \(c = \mu(B(0,r)) > 0\) by hypothesis. The balls \(B(u_n, r)\) are disjoint Borel sets, so countable additivity of \(\mu\) and Step 4 together with the monotonicity of measures give

\begin{equation*} \mu\bigl(B(0, 1+r)\bigr) \ \ge \ \mu\Bigl(\bigcup_{n=1}^{\infty} B(u_n, r)\Bigr) \ = \ \sum_{n=1}^{\infty} \mu\bigl(B(u_n, r)\bigr) \ = \ \sum_{n=1}^{\infty} c \ = \ \infty , \end{equation*}

where the last equality holds because \(c > 0\). This contradicts the hypothesis that \(\mu(B(0, 1+r)) < \infty\).

Therefore no such measure \(\mu\) exists.

Note where infinite-dimensionality was used: in a finite-dimensional Hilbert space the closed unit ball is compact, so no infinite family of disjoint balls of a fixed radius can fit inside a bounded set, and indeed Lebesgue measure (transported to \(V\) via an orthonormal basis) is a translation invariant measure assigning positive finite measure to every open ball.

Problem (8C.16)

Find the polynomial \(g\) of degree at most \(4\) that minimizes

\begin{equation*} \int_0^1 \bigl| x^5 - g(x) \bigr|^2 \, dx . \end{equation*}

Solution

Work in the real Hilbert space \(L^2([0,1])\) with the usual inner product \(\langle p, q \rangle = \int_0^1 p q\), and let

\begin{equation*} U = \operatorname{span}\{1, x, x^2, x^3, x^4\} \end{equation*}

be the subspace of polynomials of degree at most \(4\). (Distinct polynomials are distinct elements of \(L^2([0,1])\), because a nonzero polynomial has only finitely many zeros and hence is nonzero almost everywhere; so \(U\) really is a \(5\)-dimensional subspace.) Let \(f(x) = x^5\). We must find the element of \(U\) closest to \(f\); because \(U\) is finite-dimensional it is a closed subspace of \(L^2([0,1])\) (every finite-dimensional subspace of a normed vector space is closed; see Exercise 9 in Section 6D), so this element exists, is unique, and equals \(P_U f\) (by 8.28, using that a subspace is convex; see also 8.34).

If one allows polynomials with complex coefficients, nothing changes: writing \(g = g_1 + i g_2\) with \(g_1, g_2\) real polynomials of degree at most \(4\) gives \(\int_0^1 |x^5 - g|^2 = \int_0^1 (x^5 - g_1)^2 + \int_0^1 g_2^2\), which is minimized only when \(g_2 = 0\). So we may and do work with real coefficients.

Characterization. By 8.37(a) and 8.37(b), the element \(g \in U\) minimizing \(\|f - g\|\) is the unique \(g \in U\) with \(f - g \in U^{\perp}\), i.e. with

\begin{equation*} \int_0^1 \bigl( x^5 - g(x) \bigr) x^j \, dx = 0 \qquad \text{for } j = 0, 1, 2, 3, 4 . \end{equation*}

Equivalently: \(x^5 - g\) is the unique monic polynomial of degree \(5\) that is orthogonal in \(L^2([0,1])\) to every polynomial of degree at most \(4\).

Producing that monic polynomial. Let

\begin{equation*} u(x) = (x^2 - x)^5 = x^5 (x-1)^5 , \end{equation*}

and define

\begin{equation*} q = \frac{5!}{10!} \, u^{(5)} . \end{equation*}

Since \(u\) is a polynomial of degree \(10\) with leading term \(x^{10}\), its fifth derivative has leading term \(\frac{10!}{5!} x^5\); hence \(q\) is monic of degree \(5\).

I claim \(q \perp x^j\) for \(j = 0, \ldots, 4\). Both \(0\) and \(1\) are zeros of \(u\) of multiplicity \(5\), so

\begin{equation*} u(0) = u^{\prime}(0) = \cdots = u^{(4)}(0) = 0 \quad\text{and}\quad u(1) = u^{\prime}(1) = \cdots = u^{(4)}(1) = 0 . \end{equation*}

Fix \(j \in \{0, 1, 2, 3, 4\}\) and integrate by parts \(j+1\) times (legitimate, as everything in sight is a polynomial):

\begin{equation*} \int_0^1 u^{(5)}(x)\, x^j \, dx = \Bigl[ \sum_{i=0}^{j} (-1)^i u^{(4-i)}(x) \, (x^j)^{(i)} \Bigr]_0^1 + (-1)^{j+1} \int_0^1 u^{(4-j)}(x) \, (x^j)^{(j+1)} \, dx . \end{equation*}

Every boundary term involves \(u^{(4-i)}\) with \(0 \le 4 - i \le 4\), hence vanishes at both \(0\) and \(1\); and \((x^j)^{(j+1)} = 0\). Therefore \(\int_0^1 u^{(5)}(x) x^j \, dx = 0\), so \(\int_0^1 q(x) x^j \, dx = 0\) for \(j = 0, \ldots, 4\), as claimed.

Computing \(q\). Expanding,

\begin{equation*} u(x) = x^{10} - 5x^9 + 10x^8 - 10x^7 + 5x^6 - x^5 , \end{equation*}

so

\begin{equation*} u^{(5)}(x) = 30240 x^5 - 75600 x^4 + 67200 x^3 - 25200 x^2 + 3600 x - 120 . \end{equation*}

Since \(\frac{5!}{10!} = \frac{1}{30240}\), dividing gives

\begin{equation*} q(x) = x^5 - \frac{5}{2} x^4 + \frac{20}{9} x^3 - \frac{5}{6} x^2 + \frac{5}{42} x - \frac{1}{252} . \end{equation*}

(This \(q\) is the monic shifted Legendre polynomial of degree \(5\) on \([0,1]\).)

Conclusion. The minimizing polynomial is \(g = f - q\), that is,

\begin{equation*} g(x) = \frac{5}{2} x^4 - \frac{20}{9} x^3 + \frac{5}{6} x^2 - \frac{5}{42} x + \frac{1}{252} . \end{equation*}

Uniqueness of \(g\) follows from the uniqueness in 8.34. One checks directly that \(\int_0^1 (x^5 - g(x)) x^j \, dx = 0\) for \(j = 0, \ldots, 4\), confirming the answer, and the minimum value attained is

\begin{equation*} \int_0^1 \bigl( x^5 - g(x) \bigr)^2 dx = \int_0^1 q^2 = \frac{1}{698544} . \end{equation*}

Problem (8C.17)

Prove that each orthonormal family in a Hilbert space can be extended to an orthonormal basis of the Hilbert space. Specifically, suppose \(\{e_j\}_{j \in \Omega}\) is an orthonormal family in a Hilbert space \(V\). Prove that there exists a set \(\Gamma\) containing \(\Omega\) and an orthonormal basis \(\{f_k\}_{k \in \Gamma}\) of \(V\) such that \(f_j = e_j\) for every \(j \in \Omega\).

Solution

Step 0: from a family to a set. Because \(\{e_j\}_{j \in \Omega}\) is an orthonormal family, the map \(j \mapsto e_j\) is injective: if \(j \neq j^{\prime}\) then \(\langle e_j, e_{j^{\prime}} \rangle = 0\) while \(\langle e_j, e_j \rangle = 1\), so \(e_j \neq e_{j^{\prime}}\). Let

\begin{equation*} E = \{e_j : j \in \Omega\} . \end{equation*}

Then \(E\) is an orthonormal subset of \(V\) (every element has norm \(1\) and distinct elements are orthogonal), and \(j \mapsto e_j\) is a bijection from \(\Omega\) onto \(E\).

Step 1: Zorn’s Lemma. Let

\begin{equation*} \mathcal{A} = \{\varnothing\} \cup \{ S \subseteq V : S \text{ is an orthonormal subset of } V \text{ and } E \subseteq S \} . \end{equation*}

(The empty set is thrown in only so that the hypothesis of 6.60 also holds for the empty chain, whose union is \(\varnothing\).) Note \(E \in \mathcal{A}\).

Suppose \(\mathcal{C} \subseteq \mathcal{A}\) is a chain and let \(L\) be the union of all the sets in \(\mathcal{C}\). If every element of \(\mathcal{C}\) is \(\varnothing\) (in particular if \(\mathcal{C} = \varnothing\)), then \(L = \varnothing \in \mathcal{A}\). Otherwise some element of \(\mathcal{C}\) contains \(E\), so \(E \subseteq L\). If \(f \in L\), then \(f\) lies in some orthonormal set belonging to \(\mathcal{C}\), so \(\|f\| = 1\). If \(f, g \in L\) with \(f \neq g\), then there exist \(\Omega^{\prime}, \Gamma^{\prime} \in \mathcal{C}\) with \(f \in \Omega^{\prime}\) and \(g \in \Gamma^{\prime}\); because \(\mathcal{C}\) is a chain, either \(\Omega^{\prime} \subseteq \Gamma^{\prime}\) or \(\Gamma^{\prime} \subseteq \Omega^{\prime}\), so \(f\) and \(g\) lie in a common orthonormal subset of \(V\) and hence \(\langle f, g \rangle = 0\). Thus \(L\) is an orthonormal subset of \(V\) containing \(E\), i.e. \(L \in \mathcal{A}\).

Hence \(\mathcal{A}\) satisfies the hypothesis of Zorn’s Lemma (6.60) (applied with the set \(V\)), so \(\mathcal{A}\) has a maximal element \(F\).

Step 2: \(F\) is an orthonormal basis. First, \(E \subseteq F\): this is automatic unless \(F = \varnothing\), and \(F = \varnothing\) is impossible when \(E \neq \varnothing\) because then \(\varnothing \subsetneq E\) with \(E \in \mathcal{A}\), contradicting maximality. In particular \(F\) is an orthonormal subset of \(V\) containing \(E\).

Let \(\mathcal{A}_0\) be the collection of all orthonormal subsets of \(V\). We claim \(F\) is a maximal element of \(\mathcal{A}_0\). Indeed, if \(F \subsetneq F^{\prime}\) for some orthonormal subset \(F^{\prime}\) of \(V\), then \(E \subseteq F \subseteq F^{\prime}\), so \(F^{\prime} \in \mathcal{A}\), contradicting maximality of \(F\) in \(\mathcal{A}\). Thus \(F\) is a maximal element of \(\mathcal{A}_0\), and therefore \(F\) is an orthonormal basis of \(V\) by 8.74.

Step 3: indexing \(F\) by a set containing \(\Omega\). We need an index set \(\Gamma \supseteq \Omega\) and a bijection from \(\Gamma\) onto \(F\) that sends \(j\) to \(e_j\) for each \(j \in \Omega\). Put

\begin{equation*} \Lambda = (F \setminus E) \times \{\Omega\} = \bigl\{ (h, \Omega) : h \in F \setminus E \bigr\} . \end{equation*}

Then \(\Lambda\) is disjoint from \(\Omega\): if \((h, \Omega) \in \Omega\), then, writing the ordered pair in the usual way as \((h, \Omega) = \{\{h\}, \{h, \Omega\}\}\), we would get the membership cycle \(\Omega \in \{h, \Omega\} \in (h,\Omega) \in \Omega\), which is impossible. (Any other device producing a set of the same cardinality as \(F \setminus E\) and disjoint from \(\Omega\) would do just as well.)

Now let \(\Gamma = \Omega \cup \Lambda\), which contains \(\Omega\), and define \(\{f_k\}_{k \in \Gamma}\) by

\begin{equation*} f_j = e_j \ \text{ for } j \in \Omega, \qquad f_{(h, \Omega)} = h \ \text{ for } (h, \Omega) \in \Lambda . \end{equation*}

This is well defined because \(\Omega\) and \(\Lambda\) are disjoint. The map \(k \mapsto f_k\) is a bijection from \(\Gamma\) onto \(E \cup (F \setminus E) = F\): it is onto by construction, and it is injective because \(j \mapsto e_j\) is injective on \(\Omega\) with image \(E\), \((h,\Omega) \mapsto h\) is injective on \(\Lambda\) with image \(F \setminus E\), and \(E \cap (F \setminus E) = \varnothing\).

Consequently \(\{f_k\}_{k \in \Gamma}\) is an orthonormal family (distinct indices give distinct, hence orthogonal, elements of the orthonormal set \(F\), and each has norm \(1\)) whose set of values is \(F\). By the definition 8.61, being an orthonormal basis means that the closure of the span of the set of values equals \(V\); that condition depends only on the set of values, and it holds for \(F\) by Step 2. Hence \(\{f_k\}_{k \in \Gamma}\) is an orthonormal basis of \(V\). And \(f_j = e_j\) for every \(j \in \Omega\), as required.

Problem (8C.18)

Prove that every vector space has a basis.

Solution

Recall the relevant definitions. A subset \(\Gamma\) of a vector space \(V\) is linearly independent if the only way to write \(0\) as a finite linear combination \(a_1 g_1 + \cdots + a_n g_n\) with \(g_1, \ldots, g_n\) distinct elements of \(\Gamma\) and \(a_1, \ldots, a_n \in \mathbf{F}\) is with \(a_1 = \cdots = a_n = 0\); the span of \(\Gamma\) is the set of all such finite linear combinations (with \(\operatorname{span} \varnothing = \{0\}\)); and \(\Gamma\) is a basis of \(V\) if \(\Gamma\) is linearly independent and \(\operatorname{span} \Gamma = V\). This is the content of 6.61; we give the proof.

Suppose \(V\) is a vector space. Let \(\mathcal{A}\) denote the collection of all linearly independent subsets of \(V\).

The hypothesis of Zorn’s Lemma holds. Suppose \(\mathcal{C} \subseteq \mathcal{A}\) is a chain, and let \(L\) be the union of all the sets in \(\mathcal{C}\). If \(\mathcal{C} = \varnothing\), then \(L = \varnothing\), which is linearly independent (vacuously), so \(L \in \mathcal{A}\). Now suppose \(\mathcal{C} \neq \varnothing\), and suppose

\begin{equation*} a_1 g_1 + \cdots + a_n g_n = 0 \end{equation*}

with \(g_1, \ldots, g_n\) distinct elements of \(L\) and \(a_1, \ldots, a_n \in \mathbf{F}\). For each \(m \in \{1, \ldots, n\}\) there is \(\Gamma_m \in \mathcal{C}\) with \(g_m \in \Gamma_m\). Because \(\mathcal{C}\) is a chain and \(\{\Gamma_1, \ldots, \Gamma_n\}\) is a finite subset of \(\mathcal{C}\), one of \(\Gamma_1, \ldots, \Gamma_n\) contains all the others (an easy induction on \(n\): for \(n \ge 2\), one of the largest of \(\Gamma_1, \ldots, \Gamma_{n-1}\) and \(\Gamma_n\) contains the other). Call it \(\Gamma\). Then \(g_1, \ldots, g_n\) are distinct elements of the linearly independent set \(\Gamma\), so \(a_1 = \cdots = a_n = 0\). Hence \(L\) is linearly independent, so \(L \in \mathcal{A}\). In other words, linear independence is a condition checked one finite subset at a time, and each finite subset of \(L\) lies in a single element of the chain.

Applying Zorn’s Lemma. By 6.60, \(\mathcal{A}\) contains a maximal element \(\Gamma\).

A maximal linearly independent set is a basis. This is 6.57; here is the argument. Suppose \(\operatorname{span} \Gamma \neq V\), and choose \(f \in V\) with \(f \notin \operatorname{span} \Gamma\). Then \(f \notin \Gamma\) (as \(\Gamma \subseteq \operatorname{span}\Gamma\)), so \(\Gamma \subsetneq \Gamma \cup \{f\}\). We claim \(\Gamma \cup \{f\}\) is linearly independent. Suppose

\begin{equation*} a f + a_1 g_1 + \cdots + a_n g_n = 0 \end{equation*}

with \(g_1, \ldots, g_n\) distinct elements of \(\Gamma\) and \(a, a_1, \ldots, a_n \in \mathbf{F}\). If \(a \neq 0\), then

\begin{equation*} f = -\frac{a_1}{a} g_1 - \cdots - \frac{a_n}{a} g_n \in \operatorname{span} \Gamma , \end{equation*}

contradicting the choice of \(f\). Hence \(a = 0\), and then linear independence of \(\Gamma\) forces \(a_1 = \cdots = a_n = 0\). Thus \(\Gamma \cup \{f\} \in \mathcal{A}\) strictly contains \(\Gamma\), contradicting the maximality of \(\Gamma\) in \(\mathcal{A}\).

Therefore \(\operatorname{span} \Gamma = V\). Since \(\Gamma\) is linearly independent, \(\Gamma\) is a basis of \(V\).

Two remarks. First, if \(V = \{0\}\) the argument still works and produces the basis \(\varnothing\). Second, the proof is nonconstructive: Zorn’s Lemma (equivalently, the Axiom of Choice) gives existence but no recipe, which is why no one can exhibit a concrete basis of, say, \(\mathbf{R}\) as a vector space over \(\mathbf{Q}\).

Problem (8C.19)

Find the polynomial \(g\) of degree at most \(4\) such that

\begin{equation*} f\bigl(\tfrac{1}{2}\bigr) = \int_0^1 f g \end{equation*}

for every polynomial \(f\) of degree at most \(4\).

Solution

Let \(U\) denote the subspace of the real Hilbert space \(L^2([0,1])\) consisting of the polynomials of degree at most \(4\), with inner product \(\langle p, q \rangle = \int_0^1 p q\). Then \(U\) is finite-dimensional, hence a closed subspace of \(L^2([0,1])\), hence itself a Hilbert space [8.22]. Define \(\varphi : U \to \mathbf{R}\) by

\begin{equation*} \varphi(f) = f\bigl(\tfrac{1}{2}\bigr). \end{equation*}

This is well defined: a nonzero polynomial has only finitely many zeros, so distinct polynomials of degree at most \(4\) determine distinct elements of \(L^2([0,1])\), and hence each element of \(U\) arises from exactly one polynomial. Then \(\varphi\) is linear, and it is bounded because \(U\) is finite-dimensional (every linear map from a finite-dimensional normed vector space is continuous; see Exercise 7 in Section 6D). So by the Riesz Representation Theorem (8.47), there is a unique \(g \in U\) with \(\varphi(f) = \langle f, g \rangle\) for all \(f \in U\); that is, \(f(\frac12) = \int_0^1 fg\) for every polynomial \(f\) of degree at most \(4\). To compute \(g\) we use the formula 8.77 from 8.76, which says \(g = \sum_k \overline{\varphi(e_k)} \, e_k\) for any orthonormal basis \(\{e_k\}\) of \(U\). Here the scalars are real, so the complex conjugation is invisible: if \(\{e_k\}_{k=0}^{4}\) is an orthonormal basis of \(U\), then

\begin{equation*} g = \sum_{k=0}^{4} \varphi(e_k) \, e_k = \sum_{k=0}^{4} e_k\bigl(\tfrac{1}{2}\bigr) e_k . \end{equation*}

An orthonormal basis of \(U\). For \(k \in \{0, 1, 2, 3, 4\}\) define (Rodrigues’ formula for the shifted Legendre polynomials)

\begin{equation*} p_k(x) = \frac{1}{k!} \, \frac{d^k}{dx^k} \bigl[ (x^2 - x)^k \bigr] . \end{equation*}

Then \(p_k\) is a polynomial of degree \(k\) with leading coefficient \(\frac{(2k)!}{(k!)^2}\), so \(\{p_0, \ldots, p_4\}\) is a basis of \(U\).

Write \(u(x) = (x^2 - x)^k = x^k(x-1)^k\). Both \(0\) and \(1\) are zeros of \(u\) of multiplicity \(k\), so \(u, u^{\prime}, \ldots, u^{(k-1)}\) all vanish at \(0\) and at \(1\). Hence for \(0 \le j < k\), integrating by parts \(j+1\) times kills every boundary term and gives

\begin{equation*} \int_0^1 u^{(k)}(x) \, x^j \, dx = (-1)^{j+1} \int_0^1 u^{(k-j-1)}(x) \, (x^j)^{(j+1)} \, dx = 0 , \end{equation*}

because \((x^j)^{(j+1)} = 0\). Thus \(\langle p_k, x^j \rangle = 0\) for \(0 \le j < k\), and therefore \(\langle p_k, p_j \rangle = 0\) whenever \(j < k\): the \(p_k\) are pairwise orthogonal.

For the norms, integrate by parts \(k\) times (again all boundary terms vanish):

\begin{equation*} \int_0^1 u^{(k)}(x) \, x^k \, dx = (-1)^k \int_0^1 u(x) \, (x^k)^{(k)} \, dx = (-1)^k k! \int_0^1 (x^2-x)^k \, dx = k! \int_0^1 x^k (1-x)^k \, dx = k! \cdot \frac{(k!)^2}{(2k+1)!} , \end{equation*}

using the beta integral \(\int_0^1 x^k(1-x)^k dx = \frac{(k!)^2}{(2k+1)!}\). Hence \(\langle p_k, x^k \rangle = \frac{(k!)^2}{(2k+1)!}\), and since \(p_k\) is orthogonal to all polynomials of degree less than \(k\) and has leading coefficient \(\frac{(2k)!}{(k!)^2}\),

\begin{equation*} \|p_k\|^2 = \Bigl\langle p_k , \frac{(2k)!}{(k!)^2} x^k \Bigr\rangle = \frac{(2k)!}{(k!)^2} \cdot \frac{(k!)^2}{(2k+1)!} = \frac{1}{2k+1} . \end{equation*}

Therefore \(e_k = \sqrt{2k+1} \, p_k\), \(k = 0, \ldots, 4\), is an orthonormal basis of \(U\) (five orthonormal vectors in a \(5\)-dimensional space).

Evaluating at \(\frac{1}{2}\). Since \((x^2-x)^k\) is unchanged by \(x \mapsto 1-x\), differentiating \(k\) times gives \(p_k(1-x) = (-1)^k p_k(x)\); hence \(p_k(\frac12) = 0\) for odd \(k\). Computing directly,

\begin{equation*} p_0(x) = 1, \quad p_2(x) = 6x^2 - 6x + 1, \quad p_4(x) = 70x^4 - 140x^3 + 90x^2 - 20x + 1 , \end{equation*}

so

\begin{equation*} p_0\bigl(\tfrac12\bigr) = 1, \qquad p_2\bigl(\tfrac12\bigr) = -\tfrac12, \qquad p_4\bigl(\tfrac12\bigr) = \tfrac38 . \end{equation*}

The answer. Because \(e_k(\frac12) e_k = (2k+1) p_k(\frac12) p_k\),

\begin{equation*} g = 1 \cdot p_0 + 5 \cdot \bigl(-\tfrac12\bigr) p_2 + 9 \cdot \tfrac38 \, p_4 = 1 - \tfrac52 (6x^2 - 6x + 1) + \tfrac{27}{8}(70x^4 - 140x^3 + 90x^2 - 20x + 1), \end{equation*}

that is,

\begin{equation*} g(x) = \frac{945}{4} x^4 - \frac{945}{2} x^3 + \frac{1155}{4} x^2 - \frac{105}{2} x + \frac{15}{8} = \frac{15}{8}\bigl( 126 x^4 - 252 x^3 + 154 x^2 - 28 x + 1 \bigr). \end{equation*}

Check. It suffices to verify \(\int_0^1 x^j g(x) \, dx = 2^{-j}\) for \(j = 0, 1, 2, 3, 4\), since those monomials span \(U\). Indeed

\begin{equation*} \int_0^1 x^j g(x)\,dx = \frac{945}{4(j+5)} - \frac{945}{2(j+4)} + \frac{1155}{4(j+3)} - \frac{105}{2(j+2)} + \frac{15}{8(j+1)} , \end{equation*}

which equals \(1, \frac12, \frac14, \frac18, \frac{1}{16}\) for \(j = 0, 1, 2, 3, 4\) respectively.

Problem (8C.20)

Suppose \(G\) is a nonempty open subset of \(\mathbf{C}\). The Bergman space \(L^2_a(G)\) is defined to be the set of analytic functions \(f : G \to \mathbf{C}\) such that

\begin{equation*} \int_G |f|^2 \, d\lambda_2 < \infty , \end{equation*}

where \(\lambda_2\) is the usual Lebesgue measure on \(\mathbf{R}^2\), which is identified with \(\mathbf{C}\). For \(f, h \in L^2_a(G)\), define \(\langle f, h \rangle\) to be \(\int_G f \overline{h} \, d\lambda_2\).

(a) Show that \(L^2_a(G)\) is a Hilbert space.

(b) Show that if \(w \in G\), then \(f \mapsto f(w)\) is a bounded linear functional on \(L^2_a(G)\).

Solution

Throughout, \(D(w, r) = \{z \in \mathbf{C} : |z - w| < r\}\) and \(\|f\| = \bigl( \int_G |f|^2 \, d\lambda_2 \bigr)^{1/2}\).

A key pointwise estimate. We first prove the following lemma, which drives both parts.

Lemma. Suppose \(f\) is analytic on \(G\), \(w \in G\), and \(r > 0\) satisfies \(\overline{D(w,r)} \subseteq G\). Then

\begin{equation*} f(w) = \frac{1}{\pi r^2} \int_{D(w,r)} f \, d\lambda_2 , \end{equation*}

and consequently, if \(\int_G |f|^2 d\lambda_2 < \infty\), then

\begin{equation*} |f(w)| \le \frac{1}{\sqrt{\pi} \, r} \, \|f\| . \end{equation*}

Proof of Lemma. For \(0 < s \le r\), Cauchy’s integral formula applied to the circle of radius \(s\) centered at \(w\) (a circle contained in \(G\), along with its inside) gives

\begin{equation*} f(w) = \frac{1}{2\pi i} \int_{|z - w| = s} \frac{f(z)}{z - w} \, dz = \frac{1}{2\pi} \int_0^{2\pi} f(w + s e^{i\theta}) \, d\theta . \end{equation*}

Multiply by \(2\pi s\) and integrate over \(s \in (0, r)\). Since \(f\) is continuous on the compact set \(\overline{D(w,r)}\) it is bounded there, so all integrals below are absolutely convergent, and the polar-coordinates change of variables gives

\begin{equation*} \int_{D(w,r)} f \, d\lambda_2 = \int_0^r \! \int_0^{2\pi} f(w + s e^{i\theta}) \, s \, d\theta \, ds = \int_0^r 2\pi s f(w) \, ds = \pi r^2 f(w) , \end{equation*}

which is the first assertion. For the second, the Cauchy-Schwarz inequality (8.11), applied in \(L^2(\lambda_2|_{D(w,r)})\) to \(|f|\) and the constant function \(1\), gives

\begin{equation*} \int_{D(w,r)} |f| \, d\lambda_2 \le \Bigl( \int_{D(w,r)} |f|^2 \, d\lambda_2 \Bigr)^{1/2} \bigl( \lambda_2(D(w,r)) \bigr)^{1/2} \le \|f\| \sqrt{\pi} \, r . \end{equation*}

Hence \(|f(w)| \le \frac{1}{\pi r^2} \cdot \sqrt{\pi} \, r \, \|f\| = \frac{\|f\|}{\sqrt{\pi} \, r}\), completing the proof of the Lemma.

(a) \(L^2_a(G)\) is an inner product space. Sums and scalar multiples of analytic functions are analytic, and if \(f, h \in L^2_a(G)\) and \(\alpha \in \mathbf{C}\) then \(|f + h|^2 \le 2|f|^2 + 2|h|^2\) and \(|\alpha f|^2 = |\alpha|^2 |f|^2\) show that \(f + h\) and \(\alpha f\) are again square integrable. So \(L^2_a(G)\) is a vector space. For \(f, h \in L^2_a(G)\) the product \(f \overline{h}\) is integrable over \(G\) by the Cauchy-Schwarz inequality (8.11), so \(\langle f, h \rangle\) is well defined; the algebraic properties required of an inner product (linearity in the first slot, conjugate symmetry, \(\langle f, f \rangle \ge 0\)) are immediate from the corresponding properties of the integral. For definiteness, suppose \(\langle f, f \rangle = \int_G |f|^2 d\lambda_2 = 0\). Then \(|f|^2 = 0\) almost everywhere on \(G\). If \(f(w) \neq 0\) for some \(w \in G\), then by continuity \(|f| > 0\) on some ball \(D(w, \varepsilon) \subseteq G\), a set of positive measure \(\pi \varepsilon^2\), a contradiction. Hence \(f = 0\). Thus \(\langle \cdot , \cdot \rangle\) is an inner product on \(L^2_a(G)\), with associated norm \(\|\cdot\|\).

Completeness. Suppose \(f_1, f_2, \ldots\) is a Cauchy sequence in \(L^2_a(G)\).

For a compact set \(K \subseteq G\), let

\begin{equation*} r_K = \min\Bigl\{ 1, \tfrac{1}{2} \operatorname{dist}(K, \mathbf{C} \setminus G) \Bigr\} > 0 \end{equation*}

(interpreting \(\operatorname{dist}(K, \varnothing) = \infty\); positivity holds because \(K\) is compact, \(\mathbf{C} \setminus G\) is closed, and they are disjoint). Then \(\overline{D(w, r_K)} \subseteq G\) for every \(w \in K\), so the Lemma applied to the analytic function \(f_n - f_m\) gives

\begin{equation*} \sup_{w \in K} |f_n(w) - f_m(w)| \le \frac{1}{\sqrt{\pi} \, r_K} \, \|f_n - f_m\| . \end{equation*}

Hence \(f_1, f_2, \ldots\) is uniformly Cauchy on each compact subset of \(G\). Therefore this sequence converges pointwise on \(G\) to a function \(f : G \to \mathbf{C}\), and the convergence is uniform on each compact subset of \(G\). Being a limit, uniform on compact sets, of analytic functions, \(f\) is analytic on \(G\) (this is the standard Weierstrass convergence theorem, proved via Morera’s theorem: the limit is continuous, and its integral over each triangle in a disk contained in \(G\) is the limit of integrals that are \(0\)).

Now let \(\varepsilon > 0\) and choose \(N\) so that \(\|f_n - f_m\| < \varepsilon\) whenever \(n, m \ge N\). Fix \(n \ge N\) and a compact \(K \subseteq G\). Because \(f_m \to f\) uniformly on \(K\) and \(\lambda_2(K) < \infty\),

\begin{equation*} \int_K |f_n - f|^2 \, d\lambda_2 = \lim_{m \to \infty} \int_K |f_n - f_m|^2 \, d\lambda_2 \le \limsup_{m \to \infty} \|f_n - f_m\|^2 \le \varepsilon^2 . \end{equation*}

Apply this with the compact sets

\begin{equation*} K_j = \Bigl\{ z \in \mathbf{C} : |z| \le j \text{ and } \operatorname{dist}(z, \mathbf{C} \setminus G) \ge \tfrac{1}{j} \Bigr\} \qquad (j \in \mathbf{Z}^+), \end{equation*}

which are closed and bounded, hence compact, are contained in \(G\), increase with \(j\), and have union \(G\) (if \(z \in G\) then \(\operatorname{dist}(z, \mathbf{C}\setminus G) > 0\) because \(G\) is open). Since \(\chi_{K_j} |f_n - f|^2\) increases pointwise to \(|f_n - f|^2\), the Monotone Convergence Theorem (3.11) gives

\begin{equation*} \int_G |f_n - f|^2 \, d\lambda_2 = \lim_{j \to \infty} \int_{K_j} |f_n - f|^2 \, d\lambda_2 \le \varepsilon^2 . \end{equation*}

In particular \(f_n - f\) is square integrable on \(G\), hence so is \(f = f_n - (f_n - f)\); thus \(f \in L^2_a(G)\). And the displayed inequality says \(\|f_n - f\| \le \varepsilon\) for all \(n \ge N\), so \(f_n \to f\) in \(L^2_a(G)\).

Therefore every Cauchy sequence in \(L^2_a(G)\) converges in \(L^2_a(G)\), so \(L^2_a(G)\) is an inner product space that is complete, i.e. a Hilbert space (8.21).

(b) Fix \(w \in G\). The map \(\varphi : L^2_a(G) \to \mathbf{C}\) defined by \(\varphi(f) = f(w)\) is clearly linear, since addition and scalar multiplication in \(L^2_a(G)\) are pointwise. Because \(G\) is open, there exists \(r > 0\) with \(\overline{D(w, r)} \subseteq G\). The Lemma then gives

\begin{equation*} |\varphi(f)| = |f(w)| \le \frac{1}{\sqrt{\pi} \, r} \, \|f\| \end{equation*}

for every \(f \in L^2_a(G)\). Hence \(\varphi\) is a bounded linear functional, with \(\|\varphi\| \le \frac{1}{\sqrt{\pi} \, r}\).

Note that this fails badly for \(L^2(\lambda_2|_G)\) itself, where point evaluation is not even well defined; analyticity is what makes the Bergman space a space of genuine functions on which evaluation is continuous.

Problem (8C.21)

Let \(\mathbf{D}\) denote the open unit disk in \(\mathbf{C}\); thus

\begin{equation*} \mathbf{D} = \{ z \in \mathbf{C} : |z| < 1 \} . \end{equation*}

(a) Find an orthonormal basis of \(L^2_a(\mathbf{D})\).

(b) Suppose \(f \in L^2_a(\mathbf{D})\) has Taylor series

\begin{equation*} f(z) = \sum_{k=0}^{\infty} a_k z^k \end{equation*}

for \(z \in \mathbf{D}\). Find a formula for \(\|f\|\) in terms of \(a_0, a_1, a_2, \ldots\).

(c) Suppose \(w \in \mathbf{D}\). By the previous exercise and the Riesz Representation Theorem (8.47 and 8.76), there exists \(\Gamma_w \in L^2_a(\mathbf{D})\) such that

\begin{equation*} f(w) = \langle f, \Gamma_w \rangle \quad \text{for all } f \in L^2_a(\mathbf{D}). \end{equation*}

Find an explicit formula for \(\Gamma_w\).

Solution

A basic integral. For \(j, k \in \{0, 1, 2, \ldots\}\) and \(0 < s \le 1\), polar coordinates give

\begin{equation*} \int_{D(0,s)} z^j \overline{z^k} \, d\lambda_2 = \int_0^s \! \int_0^{2\pi} t^{j+k} e^{i(j-k)\theta} \, t \, d\theta \, dt = \begin{cases} 0 & \text{if } j \neq k, \\[2pt] \dfrac{\pi s^{2k+2}}{k+1} & \text{if } j = k, \end{cases} \end{equation*}

since \(\int_0^{2\pi} e^{i(j-k)\theta} d\theta = 0\) for \(j \neq k\) and \(2\pi \int_0^s t^{2k+1} dt = \frac{\pi s^{2k+2}}{k+1}\).

Taylor coefficients as inner products. Suppose \(f \in L^2_a(\mathbf{D})\) with \(f(z) = \sum_{n=0}^{\infty} a_n z^n\) for \(z \in \mathbf{D}\) (the Taylor series of \(f\) about \(0\) converges on all of \(\mathbf{D}\) because \(f\) is analytic there, and converges uniformly on each closed disk \(\overline{D(0,s)}\) with \(s < 1\)). Fix \(k \ge 0\) and \(0 < s < 1\). Uniform convergence on \(\overline{D(0,s)}\), a set of finite measure, permits term-by-term integration, so by the basic integral

\begin{equation*} \int_{D(0,s)} f(z) \overline{z^k} \, d\lambda_2 = \sum_{n=0}^{\infty} a_n \int_{D(0,s)} z^n \overline{z^k} \, d\lambda_2 = \frac{\pi s^{2k+2}}{k+1} a_k . \end{equation*}

Now let \(s \uparrow 1\). The integrands \(\chi_{D(0,s)} f \overline{z^k}\) converge pointwise on \(\mathbf{D}\) to \(f \overline{z^k}\) and are dominated by \(|f|\), which is integrable on \(\mathbf{D}\) by the Cauchy-Schwarz inequality (8.11), since \(\int_{\mathbf{D}} |f| \le \|f\| \sqrt{\pi} < \infty\). Hence the Dominated Convergence Theorem (3.31) gives

\begin{equation*} \langle f, z^k \rangle = \int_{\mathbf{D}} f \overline{z^k} \, d\lambda_2 = \frac{\pi}{k+1} a_k \qquad \text{for } k = 0, 1, 2, \ldots \end{equation*}

(a) For \(k = 0, 1, 2, \ldots\) define

\begin{equation*} e_k(z) = \sqrt{\frac{k+1}{\pi}} \, z^k . \end{equation*}

Each \(e_k\) is analytic on \(\mathbf{D}\) and, by the basic integral with \(s = 1\), satisfies \(\|z^k\|^2 = \frac{\pi}{k+1} < \infty\); thus \(e_k \in L^2_a(\mathbf{D})\) with \(\|e_k\| = 1\), and \(\langle e_j, e_k \rangle = 0\) for \(j \neq k\). So \(\{e_k\}_{k \ge 0}\) is an orthonormal family, and the \(e_k\) are distinct, so \(\{e_k : k \ge 0\}\) is an orthonormal subset of \(L^2_a(\mathbf{D})\).

It is an orthonormal basis. By 8.74 it suffices to show that this orthonormal set is a maximal orthonormal subset of \(L^2_a(\mathbf{D})\). If it were not maximal, there would be \(g \in L^2_a(\mathbf{D})\) with \(\|g\| = 1\) and \(\langle g, e_k \rangle = 0\) for every \(k\). Writing \(g(z) = \sum_k a_k z^k\), the displayed computation above gives \(\frac{\pi}{k+1} a_k = \langle g, z^k \rangle = 0\), so \(a_k = 0\) for every \(k\), so \(g = 0\), contradicting \(\|g\| = 1\). Hence

\begin{equation*} \Bigl\{ \sqrt{\tfrac{k+1}{\pi}} \, z^k \Bigr\}_{k = 0, 1, 2, \ldots} \end{equation*}

is an orthonormal basis of \(L^2_a(\mathbf{D})\).

(b) With \(f(z) = \sum_k a_k z^k\), the computation above gives

\begin{equation*} \langle f, e_k \rangle = \sqrt{\frac{k+1}{\pi}} \, \langle f, z^k \rangle = \sqrt{\frac{k+1}{\pi}} \cdot \frac{\pi a_k}{k+1} = a_k \sqrt{\frac{\pi}{k+1}} . \end{equation*}

Hence Parseval’s identity [8.63(c)] yields

\begin{equation*} \|f\|^2 = \sum_{k=0}^{\infty} |\langle f, e_k \rangle|^2 = \pi \sum_{k=0}^{\infty} \frac{|a_k|^2}{k+1} , \end{equation*}

that is,

\begin{equation*} \|f\| = \sqrt{\pi} \Bigl( \sum_{k=0}^{\infty} \frac{|a_k|^2}{k+1} \Bigr)^{1/2} . \end{equation*}

(In particular, a function \(f\) analytic on \(\mathbf{D}\) with Taylor coefficients \(a_0, a_1, \ldots\) belongs to \(L^2_a(\mathbf{D})\) if and only if \(\sum_{k} \frac{|a_k|^2}{k+1} < \infty\). One direction is the identity just proved. Conversely, if that sum is finite, then \(\sum_k a_k \sqrt{\pi/(k+1)} \, e_k\) converges in \(L^2_a(\mathbf{D})\) by 8.54, say to \(g\); the coefficient computation above applied to \(g\) shows that the Taylor coefficients of \(g\) are exactly \(a_0, a_1, \ldots\), so \(f\) and \(g\) are analytic on \(\mathbf{D}\) with the same Taylor series about \(0\), whence \(f = g \in L^2_a(\mathbf{D})\).)

(c) Let \(\varphi(f) = f(w)\), a bounded linear functional on \(L^2_a(\mathbf{D})\) by 8C.20(b). By 8.76 and 8.77, the representing vector is

\begin{equation*} \Gamma_w = \sum_{k=0}^{\infty} \overline{\varphi(e_k)} \, e_k = \sum_{k=0}^{\infty} \sqrt{\frac{k+1}{\pi}} \, \overline{w}^{\,k} \, e_k , \end{equation*}

the series converging in the norm of \(L^2_a(\mathbf{D})\). Its \(n\)-th partial sum is the polynomial

\begin{equation*} S_n(z) = \sum_{k=0}^{n} \frac{k+1}{\pi} \, \overline{w}^{\,k} z^k . \end{equation*}

Because \(|\overline{w} z| \le |w| < 1\) for \(z \in \mathbf{D}\), the identity \(\sum_{k=0}^{\infty} (k+1) u^k = \frac{1}{(1-u)^2}\) for \(|u| < 1\) shows that \(S_n\) converges pointwise on \(\mathbf{D}\) (indeed uniformly on compact subsets) to

\begin{equation*} z \mapsto \frac{1}{\pi (1 - \overline{w} z)^2} . \end{equation*}

On the other hand \(\|S_n - \Gamma_w\| \to 0\), and norm convergence in \(L^2_a(\mathbf{D})\) implies pointwise convergence, by the Lemma proved in 8C.20: for \(z \in \mathbf{D}\) and \(r = \frac{1 - |z|}{2}\) we have \(|S_n(z) - \Gamma_w(z)| \le \frac{1}{\sqrt{\pi} \, r} \|S_n - \Gamma_w\| \to 0\). The two limits must agree, so

\begin{equation*} \Gamma_w(z) = \frac{1}{\pi (1 - \overline{w} z)^2} \qquad (z \in \mathbf{D}). \end{equation*}

Check. This \(\Gamma_w\) has Taylor coefficients \(a_k = \frac{k+1}{\pi} \overline{w}^{\,k}\), so by part (b)

\begin{equation*} \|\Gamma_w\|^2 = \pi \sum_{k=0}^{\infty} \frac{1}{k+1} \cdot \frac{(k+1)^2 |w|^{2k}}{\pi^2} = \frac{1}{\pi} \sum_{k=0}^{\infty} (k+1) |w|^{2k} = \frac{1}{\pi (1 - |w|^2)^2} < \infty , \end{equation*}

confirming \(\Gamma_w \in L^2_a(\mathbf{D})\); and the last part of 8.76 says \(\|\varphi\| = \|\Gamma_w\| = \frac{1}{\sqrt{\pi} \, (1 - |w|^2)}\), which is consistent with the bound \(\|\varphi\| \le \frac{1}{\sqrt{\pi} \, r}\) from 8C.20(b), valid for every \(r < 1 - |w|\) and hence giving \(\|\varphi\| \le \frac{1}{\sqrt{\pi}\,(1-|w|)}\). Spelled out, the representation reads

\begin{equation*} f(w) = \frac{1}{\pi} \int_{\mathbf{D}} \frac{f(z)}{(1 - w \overline{z})^2} \, d\lambda_2(z) \qquad \text{for all } f \in L^2_a(\mathbf{D}), \end{equation*}

the classical Bergman kernel reproducing formula.

Problem (8C.22)

Suppose \(G\) is the annulus defined by

\begin{equation*} G = \{z \in \mathbb{C} : 1 < |z| < 2\}. \end{equation*}

(a) Find an orthonormal basis of \(L^2_a(G)\).

(b) Suppose \(f \in L^2_a(G)\) has Laurent series

\begin{equation*} f(z) = \sum_{k=-\infty}^{\infty} a_k z^k \end{equation*}

for \(z \in G\). Find a formula for \(\|f\|\) in terms of \(\ldots, a_{-1}, a_0, a_1, \ldots\).

Solution

Throughout we use polar coordinates on \(\mathbb{C} = \mathbb{R}^2\): if \(A = \{z : r_1 < |z| < r_2\}\) and \(h \colon A \to [0, \infty]\) is Borel measurable, then

\begin{equation*} \int_A h \, d\lambda_2 = \int_{r_1}^{r_2} \!\! \int_0^{2\pi} h(\rho e^{i\theta}) \, \rho \, d\theta \, d\rho , \end{equation*}

with both sides possibly \(\infty\); the same formula holds for \(h \in L^1(\lambda_2)\). (This is the usual change of variables \(d\lambda_2 = \rho \, d\rho \, d\theta\) together with Tonelli’s Theorem (5.28) in the nonnegative case and Fubini’s Theorem (5.32) in the \(L^1\) case.)

Step 1: a Laurent-coefficient formula.

Suppose \(A = \{z : r_1 < |z| < r_2\}\) is an annulus (with \(0 \le r_1 < r_2 \le \infty\)) and \(g\) is analytic on \(A\), with Laurent expansion \(g(z) = \sum_{k \in \mathbb{Z}} c_k z^k\) for \(z \in A\). The Laurent series converges uniformly on each compact subset of \(A\), in particular on each circle \(\{|z| = \rho\}\) with \(r_1 < \rho < r_2\). Because \(\{\theta \mapsto e^{ik\theta}/\sqrt{2\pi}\}_{k \in \mathbb{Z}}\) is an orthonormal family in \(L^2\big((0, 2\pi]\big)\) and \(\theta \mapsto g(\rho e^{i\theta}) = \sum_k c_k \rho^k e^{ik\theta}\) converges uniformly (hence in \(L^2\)), we get for each \(j \in \mathbb{Z}\)

\begin{equation*} \int_0^{2\pi} g(\rho e^{i\theta}) \, e^{-ij\theta} \, d\theta = 2\pi c_j \rho^{\,j} \end{equation*}

(term-by-term integration is legitimate because the convergence is uniform on the circle, which has finite measure). We shall also use repeatedly the following consequence of 8.54: if \(\{v_k\}_{k \in \Gamma}\) is an orthonormal family in a Hilbert space and \(\sum_{k \in \Gamma} |\beta_k|^2 < \infty\), then the unordered sum \(\sum_{k \in \Gamma} \beta_k v_k\) converges [8.54(a)] and

\begin{equation*} \Big\| \sum_{k \in \Gamma} \beta_k v_k \Big\|^2 = \sum_{k \in \Gamma} |\beta_k|^2 \end{equation*}

by 8.54(b). Note that this uses only orthonormality of the family, not that it be an orthonormal basis; in particular we never have to know that the trigonometric family is an orthonormal basis of \(L^2\big((0,2\pi]\big)\).

Apply this with \(v_k(\theta) = e^{ik\theta}/\sqrt{2\pi}\) and \(\beta_k = \sqrt{2\pi} \, c_k \rho^{\,k}\). Here \(\sum_k |\beta_k|^2 < \infty\): the Laurent series converges absolutely at each point of \(A\), so \(\sum_{k \in \mathbb{Z}} |c_k| \rho^{\,k} < \infty\), whence \(|c_k| \rho^{\,k}\) is bounded and \(\sum_k |c_k|^2 \rho^{2k} \le \big( \sup_k |c_k| \rho^{\,k} \big) \sum_k |c_k| \rho^{\,k} < \infty\). The resulting unordered sum agrees with the uniform (hence \(L^2\)) limit of the symmetric partial sums, namely \(\theta \mapsto g(\rho e^{i\theta})\). Therefore

\begin{equation*} \int_0^{2\pi} |g(\rho e^{i\theta})|^2 \, d\theta = 2\pi \sum_{k \in \mathbb{Z}} |c_k|^2 \rho^{2k} . \end{equation*}

Combining the last equation with the polar formula and Tonelli’s Theorem (5.28) (all terms are nonnegative) gives, as an identity in \([0, \infty]\),

\begin{equation*} \int_A |g|^2 \, d\lambda_2 = 2\pi \sum_{k \in \mathbb{Z}} |c_k|^2 \int_{r_1}^{r_2} \rho^{2k+1} \, d\rho . \tag{\(\ast\)} \end{equation*}

Step 2: the functions \(z^k\) and their norms.

Now take \(A = G\), so \(r_1 = 1\) and \(r_2 = 2\). For \(k \in \mathbb{Z}\) let \(u_k(z) = z^k\); each \(u_k\) is analytic on \(G\) and, by \((\ast)\) applied to \(g = u_k\),

\begin{equation*} \|u_k\|^2 = 2\pi \int_1^2 \rho^{2k+1} \, d\rho = \begin{cases} \dfrac{\pi \, (4^{\,k+1} - 1)}{k+1} & \text{if } k \neq -1, \\[2ex] 2\pi \ln 2 & \text{if } k = -1 . \end{cases} \end{equation*}

(For \(k \neq -1\) the quantity \((4^{k+1}-1)/(k+1)\) is positive: numerator and denominator have the same sign.) In particular \(\|u_k\| < \infty\), so \(u_k \in L^2_a(G)\) for every \(k \in \mathbb{Z}\).

For \(j, k \in \mathbb{Z}\) the product \(u_j \overline{u_k}\) is bounded on \(G\), hence in \(L^1(\lambda_2|_G)\), and the polar formula gives

\begin{equation*} \langle u_j, u_k \rangle = \int_1^2 \!\! \int_0^{2\pi} \rho^{\,j+k+1} e^{i(j-k)\theta} \, d\theta \, d\rho = 0 \qquad \text{when } j \neq k, \end{equation*}

since \(\int_0^{2\pi} e^{i(j-k)\theta} d\theta = 0\) for \(j \neq k\). Thus \(\{u_k\}_{k \in \mathbb{Z}}\) is an orthogonal family in \(L^2_a(G)\). Define

\begin{equation*} e_k = \frac{u_k}{\|u_k\|}, \qquad \text{that is,} \qquad e_k(z) = \sqrt{\frac{k+1}{\pi \, (4^{\,k+1}-1)}} \; z^k \ \ (k \neq -1), \qquad e_{-1}(z) = \frac{1}{\sqrt{2\pi \ln 2}} \cdot \frac{1}{z} . \end{equation*}

Then \(\{e_k\}_{k \in \mathbb{Z}}\) is an orthonormal family in \(L^2_a(G)\).

Step 3: \(\{e_k\}_{k \in \mathbb{Z}}\) is an orthonormal basis.

By Exercise 20 of this section, \(L^2_a(G)\) is a Hilbert space. By 8.74, an orthonormal subset of a Hilbert space is an orthonormal basis if and only if it is maximal among orthonormal subsets; and an orthonormal family is maximal exactly when the only element orthogonal to every member of the family is \(0\). So suppose \(f \in L^2_a(G)\) satisfies \(\langle f, e_k \rangle = 0\) for every \(k \in \mathbb{Z}\); equivalently \(\langle f, u_k \rangle = 0\) for every \(k\). Write the Laurent expansion \(f(z) = \sum_{m \in \mathbb{Z}} a_m z^m\) on \(G\).

Fix \(k \in \mathbb{Z}\). Since \(f, u_k \in L^2(\lambda_2|_G)\), the Cauchy–Schwarz inequality gives \(f \overline{u_k} \in L^1(\lambda_2|_G)\), so the polar formula applies:

\begin{equation*} \langle f, u_k \rangle = \int_1^2 \rho^{\,k+1} \left( \int_0^{2\pi} f(\rho e^{i\theta}) \, e^{-ik\theta} \, d\theta \right) d\rho = \int_1^2 \rho^{\,k+1} \cdot 2\pi a_k \rho^{\,k} \, d\rho = 2\pi a_k \int_1^2 \rho^{2k+1} \, d\rho = a_k \|u_k\|^2 , \end{equation*}

where the second equality is the first displayed formula of Step 1. Hence \(\langle f, u_k \rangle = 0\) forces \(a_k = 0\). As this holds for every \(k \in \mathbb{Z}\), all Laurent coefficients of \(f\) vanish, so \(f = 0\).

Therefore \(\{e_k\}_{k \in \mathbb{Z}}\) is a maximal orthonormal subset of \(L^2_a(G)\), and by 8.74 it is an orthonormal basis of \(L^2_a(G)\). This answers (a).

Step 4: the norm formula.

Let \(f \in L^2_a(G)\) have Laurent series \(f(z) = \sum_{k \in \mathbb{Z}} a_k z^k\). The computation in Step 3 gives \(\langle f, e_k \rangle = \langle f, u_k \rangle / \|u_k\| = a_k \|u_k\|\). Hence Parseval’s identity [8.63(c)] yields

\begin{equation*} \|f\|^2 = \sum_{k \in \mathbb{Z}} |\langle f, e_k \rangle|^2 = \sum_{k \in \mathbb{Z}} |a_k|^2 \, \|u_k\|^2 , \end{equation*}

that is,

\begin{equation*} \|f\| = \left( 2\pi (\ln 2) \, |a_{-1}|^2 \; + \; \pi \sum_{\substack{k \in \mathbb{Z} \\ k \neq -1}} \frac{4^{\,k+1} - 1}{k+1} \, |a_k|^2 \right)^{\!1/2} . \end{equation*}

(This also follows directly from \((\ast)\), which shows moreover that an analytic \(f\) on \(G\) with Laurent coefficients \(a_k\) belongs to \(L^2_a(G)\) if and only if the sum above is finite.)

Problem (8C.23)

Prove that if \(f \in L^2_a(\mathbb{D} \setminus \{0\})\), then \(f\) has a removable singularity at \(0\) (meaning that \(f\) can be extended to a function that is analytic on \(\mathbb{D}\)).

Solution

Here \(\mathbb{D} = \{z \in \mathbb{C} : |z| < 1\}\), so \(\mathbb{D} \setminus \{0\}\) is the punctured disk, which is the annulus \(\{z : 0 < |z| < 1\}\).

Suppose \(f \in L^2_a(\mathbb{D} \setminus \{0\})\). Because \(f\) is analytic on this annulus, it has a Laurent expansion

\begin{equation*} f(z) = \sum_{k = -\infty}^{\infty} a_k z^k \qquad \text{for } 0 < |z| < 1, \end{equation*}

the series converging uniformly on each compact subset of \(\mathbb{D} \setminus \{0\}\), in particular on each circle \(\{|z| = \rho\}\) with \(0 < \rho < 1\).

Fix \(\rho \in (0,1)\). The uniformly convergent series \(\sum_k a_k \rho^k e^{ik\theta}\) converges to \(\theta \mapsto f(\rho e^{i\theta})\) in the norm of \(L^2\big((0, 2\pi]\big)\), and \(\{\theta \mapsto e^{ik\theta}/\sqrt{2\pi}\}_{k \in \mathbb{Z}}\) is an orthonormal family in \(L^2\big((0,2\pi]\big)\). We use the following consequence of 8.54: if \(\{v_k\}\) is an orthonormal family in a Hilbert space and \(\sum_k |\beta_k|^2 < \infty\), then the unordered sum \(\sum_k \beta_k v_k\) converges [8.54(a)] and \(\big\| \sum_k \beta_k v_k \big\|^2 = \sum_k |\beta_k|^2\) [8.54(b)]. Only orthonormality is used here, not that the family be an orthonormal basis. Take \(v_k(\theta) = e^{ik\theta}/\sqrt{2\pi}\) and \(\beta_k = \sqrt{2\pi} \, a_k \rho^{\,k}\); the Laurent series converges absolutely at each point of the punctured disk, so \(\sum_{k \in \mathbb{Z}} |a_k| \rho^{\,k} < \infty\), hence \(|a_k| \rho^{\,k}\) is bounded and \(\sum_k |\beta_k|^2 = 2\pi \sum_k |a_k|^2 \rho^{2k} < \infty\). The unordered sum \(\sum_k \beta_k v_k\) then agrees with the uniform (hence \(L^2\)) limit of the symmetric partial sums, which is \(\theta \mapsto f(\rho e^{i\theta})\). Therefore

\begin{equation*} \int_0^{2\pi} |f(\rho e^{i\theta})|^2 \, d\theta = 2\pi \sum_{k = -\infty}^{\infty} |a_k|^2 \rho^{2k} . \end{equation*}

Now integrate in \(\rho\). Writing the integral over the punctured disk in polar coordinates (\(d\lambda_2 = \rho \, d\rho \, d\theta\)) and using Tonelli’s Theorem (5.28) twice — once to justify the polar formula for the nonnegative function \(|f|^2\), and once to interchange the sum over \(k \in \mathbb{Z}\) with the integral over \(\rho\), all terms being nonnegative — we obtain the identity in \([0, \infty]\)

\begin{equation*} \int_{\mathbb{D} \setminus \{0\}} |f|^2 \, d\lambda_2 = \int_0^1 \rho \left( \int_0^{2\pi} |f(\rho e^{i\theta})|^2 \, d\theta \right) d\rho = 2\pi \sum_{k = -\infty}^{\infty} |a_k|^2 \int_0^1 \rho^{2k+1} \, d\rho . \end{equation*}

By hypothesis the left side is finite, so every term on the right is finite. But if \(k \le -1\) then \(2k + 1 \le -1\), so

\begin{equation*} \int_0^1 \rho^{2k+1} \, d\rho = \infty . \end{equation*}

In \([0, \infty]\) the term indexed by such a \(k\) equals \(2\pi |a_k|^2 \cdot \infty\), which is \(\infty\) if \(a_k \neq 0\) and \(0\) if \(a_k = 0\) (using the book’s convention that \(0 \cdot \infty = 0\)). Since the whole sum is finite, \(a_k = 0\) for every \(k \le -1\).

Consequently the principal part of the Laurent series of \(f\) at \(0\) vanishes, and

\begin{equation*} f(z) = \sum_{k=0}^{\infty} a_k z^k \qquad \text{for } 0 < |z| < 1 . \end{equation*}

The right side is the analytic part of the Laurent expansion, hence a power series that converges on all of \(\mathbb{D}\); call its sum \(\tilde{f}\). Then \(\tilde{f}\) is analytic on \(\mathbb{D}\) (a convergent power series is analytic on its disk of convergence) and \(\tilde{f} = f\) on \(\mathbb{D} \setminus \{0\}\). Thus \(\tilde{f}\) is an analytic extension of \(f\) to \(\mathbb{D}\), which is exactly the statement that \(f\) has a removable singularity at \(0\).

(Equivalently, one can avoid mentioning removability directly: setting \(f(0) = a_0\) makes \(f\) analytic on \(\mathbb{D}\). Note also that the same computation shows \(\tilde{f} \in L^2_a(\mathbb{D})\) with the same norm, since \(\{0\}\) has \(\lambda_2\)-measure \(0\); so \(L^2_a(\mathbb{D} \setminus \{0\})\) and \(L^2_a(\mathbb{D})\) are naturally the same Hilbert space.)

Problem (8C.24)

The Dirichlet space \(\mathcal{D}\) is defined to be the set of analytic functions \(f \colon \mathbb{D} \to \mathbb{C}\) such that

\begin{equation*} \int_{\mathbb{D}} |f^{\prime}|^2 \, d\lambda_2 < \infty . \end{equation*}

For \(f, g \in \mathcal{D}\), define \(\langle f, g \rangle\) to be \(f(0) \overline{g(0)} + \int_{\mathbb{D}} f^{\prime} \, \overline{g^{\prime}} \, d\lambda_2\).

(a) Show that \(\mathcal{D}\) is a Hilbert space.

(b) Show that if \(w \in \mathbb{D}\), then \(f \mapsto f(w)\) is a bounded linear functional on \(\mathcal{D}\).

(c) Find an orthonormal basis of \(\mathcal{D}\).

(d) Suppose \(f \in \mathcal{D}\) has Taylor series

\begin{equation*} f(z) = \sum_{k=0}^{\infty} a_k z^k \end{equation*}

for \(z \in \mathbb{D}\). Find a formula for \(\|f\|\) in terms of \(a_0, a_1, a_2, \ldots\).

(e) Suppose \(w \in \mathbb{D}\). Find an explicit formula for \(\Gamma_w \in \mathcal{D}\) such that

\begin{equation*} f(w) = \langle f, \Gamma_w \rangle \quad \text{for all } f \in \mathcal{D} . \end{equation*}

Solution

Key computation. Let \(f\) be any analytic function on \(\mathbb{D}\), with Taylor series \(f(z) = \sum_{k=0}^{\infty} a_k z^k\) (which converges on \(\mathbb{D}\)). Then \(f^{\prime}(z) = \sum_{k=1}^{\infty} k a_k z^{k-1}\), with the same disk of convergence, and I claim that as an identity in \([0, \infty]\)

\begin{equation*} \int_{\mathbb{D}} |f^{\prime}|^2 \, d\lambda_2 = \pi \sum_{k=1}^{\infty} k \, |a_k|^2 . \tag{\(\dagger\)} \end{equation*}

Indeed, fix \(\rho \in (0,1)\). The series \(\sum_{k \ge 1} k a_k \rho^{k-1} e^{i(k-1)\theta}\) converges uniformly in \(\theta\) (a power series converges uniformly on compact subsets of its disk of convergence), hence converges to \(\theta \mapsto f^{\prime}(\rho e^{i\theta})\) in \(L^2\big((0,2\pi]\big)\). The family \(\{\theta \mapsto e^{ij\theta}/\sqrt{2\pi}\}_{j \in \mathbb{Z}}\) is orthonormal in \(L^2\big((0,2\pi]\big)\), and \(\sum_{k \ge 1} \big| \sqrt{2\pi} \, k a_k \rho^{k-1} \big|^2 < \infty\) (the power series for \(f^{\prime}\) converges absolutely at \(\rho\), so \(k|a_k|\rho^{k-1}\) is a bounded summable family). Hence by 8.54(a) the unordered sum converges, it agrees with the uniform limit \(\theta \mapsto f^{\prime}(\rho e^{i\theta})\), and 8.54(b) gives

\begin{equation*} \int_0^{2\pi} |f^{\prime}(\rho e^{i\theta})|^2 \, d\theta = 2\pi \sum_{k=1}^{\infty} k^2 |a_k|^2 \rho^{2k-2} . \end{equation*}

Writing the integral over \(\mathbb{D}\) in polar coordinates (\(d\lambda_2 = \rho \, d\rho \, d\theta\)) and interchanging the sum and the \(\rho\)-integral, both steps being legitimate for nonnegative integrands by Tonelli’s Theorem (5.28),

\begin{equation*} \int_{\mathbb{D}} |f^{\prime}|^2 \, d\lambda_2 = 2\pi \sum_{k=1}^{\infty} k^2 |a_k|^2 \int_0^1 \rho^{2k-1} \, d\rho = 2\pi \sum_{k=1}^{\infty} k^2 |a_k|^2 \cdot \frac{1}{2k} = \pi \sum_{k=1}^{\infty} k |a_k|^2 , \end{equation*}

which is \((\dagger)\). In particular an analytic \(f\) on \(\mathbb{D}\) lies in \(\mathcal{D}\) if and only if \(\sum_{k \ge 1} k |a_k|^2 < \infty\), and then

\begin{equation*} \|f\|^2 = |a_0|^2 + \pi \sum_{k=1}^{\infty} k \, |a_k|^2 . \tag{\(\ddagger\)} \end{equation*}

(d) Formula \((\ddagger)\) answers (d): if \(f \in \mathcal{D}\) has Taylor coefficients \(a_0, a_1, a_2, \ldots\), then

\begin{equation*} \|f\| = \left( |a_0|^2 + \pi \sum_{k=1}^{\infty} k \, |a_k|^2 \right)^{\!1/2} . \end{equation*}

(a) First, \(\mathcal{D}\) is a vector space: if \(f, g\) are analytic on \(\mathbb{D}\) and \(c \in \mathbb{C}\), then \((f + cg)^{\prime} = f^{\prime} + c g^{\prime}\), and the triangle inequality in \(L^2(\lambda_2|_{\mathbb{D}})\) gives \(\|f^{\prime} + cg^{\prime}\|_{L^2} \le \|f^{\prime}\|_{L^2} + |c| \, \|g^{\prime}\|_{L^2} < \infty\) when \(f, g \in \mathcal{D}\).

Second, \(\langle \cdot, \cdot \rangle\) is a well-defined inner product on \(\mathcal{D}\). It is well defined and finite because \(f^{\prime}, g^{\prime} \in L^2(\lambda_2|_{\mathbb{D}})\) implies \(f^{\prime} \overline{g^{\prime}} \in L^1(\lambda_2|_{\mathbb{D}})\) by the Cauchy–Schwarz inequality. It is linear in the first slot and conjugate symmetric, because both \((f, g) \mapsto f(0)\overline{g(0)}\) and \((f, g) \mapsto \int_{\mathbb{D}} f^{\prime} \overline{g^{\prime}} \, d\lambda_2\) are. Finally \(\langle f, f \rangle = |f(0)|^2 + \int_{\mathbb{D}} |f^{\prime}|^2 \, d\lambda_2 \ge 0\), and if \(\langle f, f \rangle = 0\) then \(\int_{\mathbb{D}} |f^{\prime}|^2 \, d\lambda_2 = 0\), so \(f^{\prime} = 0\) almost everywhere on \(\mathbb{D}\); as \(f^{\prime}\) is continuous, \(f^{\prime} \equiv 0\) on \(\mathbb{D}\), so \(f\) is constant on the connected open set \(\mathbb{D}\), and \(f(0) = 0\) then forces \(f = 0\).

Third, completeness. Define \(T \colon \mathcal{D} \to \ell^2\) by

\begin{equation*} T f = \big( a_0, \sqrt{\pi \cdot 1} \, a_1, \sqrt{\pi \cdot 2} \, a_2, \sqrt{\pi \cdot 3} \, a_3, \ldots \big), \end{equation*}

where \(a_0, a_1, a_2, \ldots\) are the Taylor coefficients of \(f\). By \((\ddagger)\), \(\|Tf\|_{\ell^2}^2 = |a_0|^2 + \pi \sum_{k \ge 1} k |a_k|^2 = \|f\|^2 < \infty\), so \(T\) indeed maps into \(\ell^2\) and is an isometry. Since the Taylor coefficients depend linearly on \(f\), the map \(T\) is linear; being an isometry, it is injective.

\(T\) is surjective: given \(c = (c_0, c_1, c_2, \ldots) \in \ell^2\), set \(a_0 = c_0\) and \(a_k = c_k / \sqrt{\pi k}\) for \(k \ge 1\). Because \((c_k)\) is a bounded sequence, say \(|c_k| \le M\), we have \(|a_k| \le M/\sqrt{\pi k}\) for \(k \ge 1\), hence \(\limsup_{k \to \infty} |a_k|^{1/k} \le 1\) and the power series \(f(z) = \sum_{k \ge 0} a_k z^k\) has radius of convergence at least \(1\); thus \(f\) is analytic on \(\mathbb{D}\). By \((\dagger)\), \(\int_{\mathbb{D}} |f^{\prime}|^2 d\lambda_2 = \pi \sum_{k \ge 1} k |a_k|^2 = \sum_{k \ge 1} |c_k|^2 < \infty\), so \(f \in \mathcal{D}\), and clearly \(Tf = c\).

So \(T\) is a linear isometry of \(\mathcal{D}\) onto \(\ell^2\). Now suppose \((f_n)\) is a Cauchy sequence in \(\mathcal{D}\). Then \((Tf_n)\) is Cauchy in \(\ell^2\), which is a Hilbert space (8.22), so \(Tf_n \to c\) for some \(c \in \ell^2\). Let \(f = T^{-1} c \in \mathcal{D}\). Then \(\|f_n - f\| = \|Tf_n - c\|_{\ell^2} \to 0\). Hence \(\mathcal{D}\) is complete, and therefore \(\mathcal{D}\) is a Hilbert space.

(b) Fix \(w \in \mathbb{D}\). Linearity of \(f \mapsto f(w)\) is clear. For boundedness, let \(f \in \mathcal{D}\) have Taylor coefficients \(a_0, a_1, \ldots\); since \(|w| < 1\), the Taylor series converges at \(w\) and

\begin{equation*} |f(w)| \le |a_0| + \sum_{k=1}^{\infty} |a_k| \, |w|^k = |a_0| + \sum_{k=1}^{\infty} \big( \sqrt{\pi k} \, |a_k| \big) \frac{|w|^k}{\sqrt{\pi k}} \le |a_0| + \left( \pi \sum_{k=1}^{\infty} k |a_k|^2 \right)^{\!1/2} \left( \frac{1}{\pi} \sum_{k=1}^{\infty} \frac{|w|^{2k}}{k} \right)^{\!1/2}, \end{equation*}

by the Cauchy–Schwarz inequality. Since \(\sum_{k \ge 1} t^k / k = \ln \frac{1}{1-t}\) for \(0 \le t < 1\), the last factor equals \(t_w := \big( \frac{1}{\pi} \ln \frac{1}{1 - |w|^2} \big)^{1/2}\), which is finite. Writing \(A = |a_0|\) and \(B = \big( \pi \sum_{k \ge 1} k |a_k|^2 \big)^{1/2}\), so that \(\|f\|^2 = A^2 + B^2\) by \((\ddagger)\), the Cauchy–Schwarz inequality in \(\mathbb{R}^2\) gives

\begin{equation*} |f(w)| \le A \cdot 1 + B \cdot t_w \le \sqrt{A^2 + B^2} \, \sqrt{1 + t_w^{\,2}} = \|f\| \left( 1 + \frac{1}{\pi} \ln \frac{1}{1 - |w|^2} \right)^{\!1/2} . \end{equation*}

Thus \(f \mapsto f(w)\) is a bounded linear functional on \(\mathcal{D}\), with norm at most \(\big( 1 + \frac{1}{\pi} \ln \frac{1}{1-|w|^2} \big)^{1/2}\) (part (e) shows this is its exact norm).

(c) Define \(e_0, e_1, e_2, \ldots \in \mathcal{D}\) by

\begin{equation*} e_0(z) = 1, \qquad e_k(z) = \frac{z^k}{\sqrt{\pi k}} \ \ \text{for } k \ge 1 . \end{equation*}

Each \(e_k\) is a polynomial, hence lies in \(\mathcal{D}\). Under the map \(T\) of part (a), \(T e_k\) is the element of \(\ell^2\) whose \(k\)th coordinate is \(1\) and whose other coordinates are \(0\) (index the coordinates by \(k = 0, 1, 2, \ldots\)). These are the standard basis vectors of \(\ell^2\), which form an orthonormal basis of \(\ell^2\) by Exercise 10(a) of this section. Because \(T\) is a linear isometry of \(\mathcal{D}\) onto \(\ell^2\), it preserves inner products (by the polarization identity, or directly: \(\langle f, g \rangle\) is determined by the norms of \(f + g\), \(f - g\), \(f + ig\), \(f - ig\)), and it carries orthonormal bases to orthonormal bases. Hence \(\{e_k\}_{k \ge 0}\) is an orthonormal basis of \(\mathcal{D}\).

(Direct verification, if preferred: \(\langle e_0, e_0 \rangle = 1 \cdot 1 + 0 = 1\); for \(k \ge 1\), \(\langle e_0, e_k \rangle = 1 \cdot \overline{0} + \int_{\mathbb{D}} 0 \cdot \overline{e_k^{\prime}} \, d\lambda_2 = 0\); and for \(j, k \ge 1\), \(\langle z^j, z^k \rangle = jk \int_{\mathbb{D}} z^{j-1} \overline{z^{k-1}} \, d\lambda_2\), which vanishes for \(j \neq k\) by the \(\theta\)-integral and equals \(k^2 \cdot 2\pi \int_0^1 \rho^{2k-1} d\rho = \pi k\) for \(j = k\). Completeness follows since \(\langle f, e_0 \rangle = a_0\) and \(\langle f, e_k \rangle = \sqrt{\pi k} \, a_k\), so \(f \perp e_k\) for all \(k\) forces every Taylor coefficient of \(f\) to vanish.)

(e) By part (b) and the Riesz Representation Theorem (8.76), the representing vector for the functional \(\varphi(f) = f(w)\) is \(\Gamma_w = \sum_{k=0}^{\infty} \overline{\varphi(e_k)} \, e_k\). Here \(\varphi(e_0) = 1\) and \(\varphi(e_k) = w^k / \sqrt{\pi k}\) for \(k \ge 1\), so

\begin{equation*} \Gamma_w(z) = 1 + \sum_{k=1}^{\infty} \frac{\overline{w}^{\,k}}{\sqrt{\pi k}} \cdot \frac{z^k}{\sqrt{\pi k}} = 1 + \frac{1}{\pi} \sum_{k=1}^{\infty} \frac{(\overline{w} z)^k}{k} . \end{equation*}

Since \(|\overline{w} z| < 1\) for \(z \in \mathbb{D}\), the number \(1 - \overline{w} z\) lies in the open right half-plane, and the series above sums to \(-\operatorname{Log}(1 - \overline{w} z)\), where \(\operatorname{Log}\) denotes the principal branch of the logarithm. Hence

\begin{equation*} \Gamma_w(z) = 1 - \frac{1}{\pi} \operatorname{Log}(1 - \overline{w} z) = 1 + \frac{1}{\pi} \operatorname{Log} \frac{1}{1 - \overline{w} z} . \end{equation*}

This function is analytic on \(\mathbb{D}\), and its Taylor coefficients \(b_0 = 1\), \(b_k = \overline{w}^{\,k} / (\pi k)\) satisfy \(\pi \sum_{k \ge 1} k |b_k|^2 = \frac{1}{\pi} \sum_{k \ge 1} |w|^{2k} / k = \frac{1}{\pi} \ln \frac{1}{1 - |w|^2} < \infty\), so \(\Gamma_w \in \mathcal{D}\) with

\begin{equation*} \|\Gamma_w\| = \left( 1 + \frac{1}{\pi} \ln \frac{1}{1 - |w|^2} \right)^{\!1/2} . \end{equation*}

As a check, if \(f \in \mathcal{D}\) has Taylor coefficients \(a_0, a_1, \ldots\), then by \((\ddagger)\)-style polarization (equivalently, by transporting the \(\ell^2\) inner product through \(T\)),

\begin{equation*} \langle f, \Gamma_w \rangle = a_0 \overline{b_0} + \pi \sum_{k=1}^{\infty} k \, a_k \overline{b_k} = a_0 + \pi \sum_{k=1}^{\infty} k \, a_k \frac{w^k}{\pi k} = \sum_{k=0}^{\infty} a_k w^k = f(w), \end{equation*}

as required. Comparing with part (b), the norm of the evaluation functional at \(w\) is exactly \(\|\Gamma_w\|\), so the bound found in (b) is sharp.

Problem (8C.25)

(a) Prove that the Dirichlet space \(\mathcal{D}\) is contained in the Bergman space \(L^2_a(\mathbb{D})\).

(b) Prove that there exists a function \(f \in L^2_a(\mathbb{D})\) such that \(f\) is uniformly continuous on \(\mathbb{D}\) and \(f \notin \mathcal{D}\).

Solution

We use two coefficient identities, each valid in \([0, \infty]\) for an arbitrary analytic function \(f\) on \(\mathbb{D}\) with Taylor series \(f(z) = \sum_{k=0}^{\infty} a_k z^k\):

\begin{equation*} \int_{\mathbb{D}} |f|^2 \, d\lambda_2 = \pi \sum_{k=0}^{\infty} \frac{|a_k|^2}{k+1} \qquad \text{and} \qquad \int_{\mathbb{D}} |f^{\prime}|^2 \, d\lambda_2 = \pi \sum_{k=1}^{\infty} k \, |a_k|^2 . \end{equation*}

The second is \((\dagger)\) from Exercise 24. The first is proved the same way: for \(\rho \in (0,1)\) the power series converges absolutely and uniformly on the circle of radius \(\rho\), hence in \(L^2\big((0,2\pi]\big)\), and \(\sum_{k \ge 0} |a_k|^2 \rho^{2k} < \infty\); so 8.54, applied to the orthonormal family \(\{\theta \mapsto e^{ik\theta}/\sqrt{2\pi}\}_{k \in \mathbb{Z}}\) with coefficients \(\sqrt{2\pi} \, a_k \rho^{\,k}\), gives \(\int_0^{2\pi} |f(\rho e^{i\theta})|^2 \, d\theta = 2\pi \sum_{k \ge 0} |a_k|^2 \rho^{2k}\); integrating in polar coordinates and interchanging sum and integral by Tonelli’s Theorem (5.28) yields

\begin{equation*} \int_{\mathbb{D}} |f|^2 \, d\lambda_2 = 2\pi \sum_{k=0}^{\infty} |a_k|^2 \int_0^1 \rho^{2k+1} \, d\rho = \pi \sum_{k=0}^{\infty} \frac{|a_k|^2}{k+1} . \end{equation*}

(This is Exercise 21(b).) Consequently, an analytic \(f\) on \(\mathbb{D}\) with Taylor coefficients \(a_0, a_1, \ldots\) satisfies

\begin{equation*} f \in L^2_a(\mathbb{D}) \iff \sum_{k=0}^{\infty} \frac{|a_k|^2}{k+1} < \infty, \qquad f \in \mathcal{D} \iff \sum_{k=1}^{\infty} k \, |a_k|^2 < \infty . \end{equation*}

(a) Suppose \(f \in \mathcal{D}\), with Taylor coefficients \(a_0, a_1, a_2, \ldots\), so that \(\sum_{k \ge 1} k |a_k|^2 < \infty\). For every \(k \ge 1\) we have \(\frac{1}{k+1} \le 1 \le k\), hence \(\frac{|a_k|^2}{k+1} \le k |a_k|^2\). Therefore

\begin{equation*} \sum_{k=0}^{\infty} \frac{|a_k|^2}{k+1} = |a_0|^2 + \sum_{k=1}^{\infty} \frac{|a_k|^2}{k+1} \le |a_0|^2 + \sum_{k=1}^{\infty} k \, |a_k|^2 < \infty . \end{equation*}

By the criterion above, \(f \in L^2_a(\mathbb{D})\). Thus \(\mathcal{D} \subseteq L^2_a(\mathbb{D})\).

Moreover the inclusion is continuous: because \(\pi \ge 1\),

\begin{equation*} \|f\|_{L^2_a(\mathbb{D})}^2 = \pi \sum_{k=0}^{\infty} \frac{|a_k|^2}{k+1} \le \pi \left( |a_0|^2 + \sum_{k=1}^{\infty} k |a_k|^2 \right) \le \pi \left( |a_0|^2 + \pi \sum_{k=1}^{\infty} k |a_k|^2 \right) = \pi \, \|f\|_{\mathcal{D}}^2 , \end{equation*}

so \(\|f\|_{L^2_a(\mathbb{D})} \le \sqrt{\pi} \, \|f\|_{\mathcal{D}}\) for all \(f \in \mathcal{D}\).

(b) Define

\begin{equation*} f(z) = \sum_{k=1}^{\infty} \frac{z^{2^k}}{2^{k/2}} \qquad (z \in \mathbb{D}) . \end{equation*}

Equivalently \(f(z) = \sum_{n=0}^{\infty} a_n z^n\) where \(a_n = 2^{-k/2}\) if \(n = 2^k\) for some integer \(k \ge 1\), and \(a_n = 0\) otherwise.

\(f\) is analytic on \(\mathbb{D}\). The coefficients satisfy \(|a_n| \le 1\) for all \(n\), so the power series has radius of convergence at least \(1\); hence \(f\) is analytic on \(\mathbb{D}\).

\(f\) is uniformly continuous on \(\mathbb{D}\). For \(z\) in the closed disk \(\overline{\mathbb{D}}\) we have \(|z^{2^k} / 2^{k/2}| \le 2^{-k/2}\), and \(\sum_{k \ge 1} 2^{-k/2} < \infty\). By the Weierstrass \(M\)-test the series converges uniformly on \(\overline{\mathbb{D}}\), so its sum is a continuous function on \(\overline{\mathbb{D}}\) extending \(f\). Since \(\overline{\mathbb{D}}\) is a compact metric space, a continuous function on it is uniformly continuous; restricting to \(\mathbb{D}\) shows \(f\) is uniformly continuous on \(\mathbb{D}\).

\(f \in L^2_a(\mathbb{D})\). From the estimate above, \(|f(z)| \le \sum_{k \ge 1} 2^{-k/2} = \frac{1}{\sqrt{2} - 1}\) for all \(z \in \mathbb{D}\), so \(f\) is bounded. Since \(\lambda_2(\mathbb{D}) = \pi < \infty\), we get \(\int_{\mathbb{D}} |f|^2 \, d\lambda_2 \le \frac{\pi}{(\sqrt{2}-1)^2} < \infty\). Hence \(f \in L^2_a(\mathbb{D})\).

\(f \notin \mathcal{D}\). By the second identity at the start of this solution,

\begin{equation*} \int_{\mathbb{D}} |f^{\prime}|^2 \, d\lambda_2 = \pi \sum_{n=1}^{\infty} n \, |a_n|^2 = \pi \sum_{k=1}^{\infty} 2^k \cdot \big(2^{-k/2}\big)^2 = \pi \sum_{k=1}^{\infty} 2^k \cdot 2^{-k} = \pi \sum_{k=1}^{\infty} 1 = \infty . \end{equation*}

Thus \(f \notin \mathcal{D}\).

So \(f\) is an element of \(L^2_a(\mathbb{D})\) that is uniformly continuous on \(\mathbb{D}\) yet does not belong to \(\mathcal{D}\). In particular, combined with (a), the inclusion \(\mathcal{D} \subseteq L^2_a(\mathbb{D})\) is strict, and membership in the Dirichlet space is not implied even by very good boundary behavior.

Real and Complex Measures

Exercises 9A

Problem (9A.1)

Prove or give a counterexample: If \(\nu\) is a real measure on a measurable space \((X, \mathcal{S})\) and \(A, B \in \mathcal{S}\) are such that \(\nu(A) \ge 0\) and \(\nu(B) \ge 0\), then \(\nu(A \cup B) \ge 0\).

Solution

The statement is false; here is a counterexample.

Let \(X = \{1, 2, 3\}\) and let \(\mathcal{S}\) be the \(\sigma\)-algebra of all subsets of \(X\). Define \(\nu : \mathcal{S} \to \mathbb{R}\) by

\begin{equation*} \nu(E) = 3\,\chi_E(1) - 2\,\chi_E(2) - 2\,\chi_E(3), \end{equation*}

in other words \(\nu(E)\) is the sum of the weights \(w_1 = 3\), \(w_2 = -2\), \(w_3 = -2\) over the points of \(E\). Explicitly,

\begin{equation*} \nu(\varnothing) = 0,\quad \nu(\{1\}) = 3,\quad \nu(\{2\}) = \nu(\{3\}) = -2, \end{equation*}

\begin{equation*} \nu(\{1,2\}) = \nu(\{1,3\}) = 1,\quad \nu(\{2,3\}) = -4,\quad \nu(X) = -1. \end{equation*}

First we check that \(\nu\) is a real measure on \((X, \mathcal{S})\). Clearly \(\nu(E) \in \mathbb{R}\) for every \(E \in \mathcal{S}\). Suppose \(E_1, E_2, \ldots\) is a disjoint sequence of sets in \(\mathcal{S}\). Because \(X\) has only three elements, at most three of the sets \(E_k\) are nonempty, and each of the points \(1, 2, 3\) lies in at most one of the sets \(E_k\). Hence

\begin{equation*} \sum_{k=1}^{\infty} \nu(E_k) = \sum_{k=1}^{\infty} \sum_{j=1}^{3} w_j\, \chi_{E_k}(j) = \sum_{j=1}^{3} w_j \sum_{k=1}^{\infty} \chi_{E_k}(j) = \sum_{j=1}^{3} w_j\, \chi_{\bigcup_k E_k}(j) = \nu\Bigl( \bigcup_{k=1}^{\infty} E_k \Bigr), \end{equation*}

where interchanging the two sums is legitimate because only finitely many terms are nonzero. Thus \(\nu\) is countably additive and therefore is a real measure.

Now let \(A = \{1, 2\}\) and \(B = \{1, 3\}\). Then

\begin{equation*} \nu(A) = 3 - 2 = 1 \ge 0 \qquad \text{and} \qquad \nu(B) = 3 - 2 = 1 \ge 0, \end{equation*}

but

\begin{equation*} \nu(A \cup B) = \nu(\{1,2,3\}) = 3 - 2 - 2 = -1 < 0. \end{equation*}

Thus the assertion fails.

The mechanism behind this counterexample is visible in 9.7(b), which gives

\begin{equation*} \nu(A \cup B) = \nu(A) + \nu(B) - \nu(A \cap B). \end{equation*}

For a real measure, \(\nu(A \cap B)\) can be strictly larger than \(\nu(A) + \nu(B)\), and then \(\nu(A \cup B) < 0\) even though \(\nu(A) \ge 0\) and \(\nu(B) \ge 0\). In the example above, \(\nu(A \cap B) = \nu(\{1\}) = 3 > 2 = \nu(A) + \nu(B)\). Of course, if \(A\) and \(B\) are disjoint then \(\nu(A \cap B) = \nu(\varnothing) = 0\) by 9.3(a), and the assertion does hold in that special case.

Problem (9A.2)

Suppose \(\nu\) is a real measure on \((X, \mathcal{S})\). Define \(\mu : \mathcal{S} \to [0, \infty)\) by

\begin{equation*} \mu(E) = |\nu(E)|. \end{equation*}

Prove that \(\mu\) is a (positive) measure on \((X, \mathcal{S})\) if and only if the range of \(\nu\) is contained in \([0, \infty)\) or the range of \(\nu\) is contained in \((-\infty, 0]\).

Solution

First suppose the range of \(\nu\) is contained in \([0, \infty)\). Then \(\mu(E) = |\nu(E)| = \nu(E)\) for every \(E \in \mathcal{S}\), so \(\mu = \nu\). Now \(\mu(\varnothing) = \nu(\varnothing) = 0\) by 9.3(a), and \(\mu\) is countably additive because \(\nu\) is. Because \(\mu\) takes values in \([0, \infty)\), we conclude that \(\mu\) is a (positive) measure.

Next suppose the range of \(\nu\) is contained in \((-\infty, 0]\). Then \(\mu(E) = |\nu(E)| = -\nu(E)\) for every \(E \in \mathcal{S}\), so \(\mu = -\nu\). Again \(\mu(\varnothing) = 0\) and \(\mu\) is countably additive (a scalar multiple of a countably additive function is countably additive), and \(\mu\) takes values in \([0, \infty)\). Thus \(\mu\) is a (positive) measure.

To prove the implication in the other direction, suppose now that \(\mu\) is a (positive) measure on \((X, \mathcal{S})\).

The first step is to record a sign condition. Suppose \(E, F \in \mathcal{S}\) are disjoint. Countable additivity of \(\mu\) applied to the disjoint sequence \(E, F, \varnothing, \varnothing, \ldots\) (note \(\mu(\varnothing) = 0\)) gives \(\mu(E \cup F) = \mu(E) + \mu(F)\), which by the definition of \(\mu\) and the additivity of \(\nu\) says

\begin{equation*} |\nu(E) + \nu(F)| = |\nu(E \cup F)| = |\nu(E)| + |\nu(F)|. \end{equation*}

Squaring both sides of \(|\nu(E) + \nu(F)| = |\nu(E)| + |\nu(F)|\) and using \(\nu(E), \nu(F) \in \mathbb{R}\) yields

\begin{equation*} \nu(E)^2 + 2\nu(E)\nu(F) + \nu(F)^2 = \nu(E)^2 + 2|\nu(E)|\,|\nu(F)| + \nu(F)^2, \end{equation*}

hence \(\nu(E)\nu(F) = |\nu(E)|\,|\nu(F)| \ge 0\). Thus

9A.2.a: \(\nu(E)\nu(F) \ge 0\) for all disjoint \(E, F \in \mathcal{S}\).

In words, \(\nu\) cannot take a strictly positive value and a strictly negative value on two disjoint sets.

Now suppose, seeking a contradiction, that the range of \(\nu\) is contained neither in \([0, \infty)\) nor in \((-\infty, 0]\). Then there exist \(A, B \in \mathcal{S}\) with

\begin{equation*} \nu(A) > 0 \qquad \text{and} \qquad \nu(B) < 0. \end{equation*}

Consider the three pairwise disjoint sets

\begin{equation*} A \cap B, \qquad A \setminus B, \qquad B \setminus A, \end{equation*}

all of which belong to \(\mathcal{S}\). Note that \(\nu\) is finitely additive: if \(C, D \in \mathcal{S}\) are disjoint, then applying countable additivity to the disjoint sequence \(C, D, \varnothing, \varnothing, \ldots\) and using \(\nu(\varnothing) = 0\) (9.3(a)) gives \(\nu(C \cup D) = \nu( C) + \nu(D)\). Since \(A\) is the disjoint union of \(A \cap B\) and \(A \setminus B\), and \(B\) is the disjoint union of \(A \cap B\) and \(B \setminus A\), we obtain

\begin{equation*} \nu(A) = \nu(A \cap B) + \nu(A \setminus B) \qquad \text{and} \qquad \nu(B) = \nu(A \cap B) + \nu(B \setminus A). \end{equation*}

I claim that the three numbers \(\nu(A \cap B)\), \(\nu(A \setminus B)\), \(\nu(B \setminus A)\) are either all \(\ge 0\) or all \(\le 0\). Indeed, if this were not so, then one of these three numbers would be strictly negative and another one (necessarily a different one of the three) would be strictly positive; because the corresponding two sets are disjoint, the product of these two values would be strictly negative, contradicting 9A.2.a.

If all three numbers are \(\ge 0\), then \(\nu(B) = \nu(A \cap B) + \nu(B \setminus A) \ge 0\), contradicting \(\nu(B) < 0\). If all three numbers are \(\le 0\), then \(\nu(A) = \nu(A \cap B) + \nu(A \setminus B) \le 0\), contradicting \(\nu(A) > 0\). Either way we have a contradiction.

Hence the range of \(\nu\) is contained in \([0, \infty)\) or is contained in \((-\infty, 0]\), completing the proof.

Problem (9A.3)

Suppose \(\nu\) is a complex measure on a measurable space \((X, \mathcal{S})\). Prove that \(|\nu|(X) = \nu(X)\) if and only if \(\nu\) is a (positive) measure.

Solution

Note that \(|\nu|(X) \in [0, \infty]\), so the equation \(|\nu|(X) = \nu(X)\) asserts in particular that the complex number \(\nu(X)\) is real and nonnegative.

First suppose \(\nu\) is a (positive) measure. Because \(\nu\) is a complex measure, \(\nu(E) \in \mathbb{C}\) for every \(E \in \mathcal{S}\); combined with \(\nu\) being a positive measure this means \(\nu(E) \in [0, \infty)\) for every \(E \in \mathcal{S}\). In particular \(\nu\) is a finite positive measure.

Taking \(n = 1\) and \(E_1 = X\) in the definition 9.8 of the total variation measure gives

\begin{equation*} |\nu|(X) \ge |\nu(X)| = \nu(X). \end{equation*}

For the reverse inequality, suppose \(n \in \mathbb{Z}^{+}\) and \(E_1, \ldots, E_n\) are disjoint sets in \(\mathcal{S}\) with \(E_1 \cup \cdots \cup E_n \subseteq X\). Because \(\nu\) is a positive measure, \(|\nu(E_k)| = \nu(E_k)\) for each \(k\), and hence

\begin{equation*} |\nu(E_1)| + \cdots + |\nu(E_n)| = \nu(E_1) + \cdots + \nu(E_n) = \nu(E_1 \cup \cdots \cup E_n) \le \nu(X), \end{equation*}

where the last inequality holds because \(\nu\) is a positive measure and hence is monotone. Taking the supremum over all such choices gives \(|\nu|(X) \le \nu(X)\). Therefore \(|\nu|(X) = \nu(X)\).

To prove the implication in the other direction, suppose now that \(|\nu|(X) = \nu(X)\). In particular \(\nu(X)\) is a nonnegative real number.

Let \(E \in \mathcal{S}\). Then \(E\) and \(X \setminus E\) are disjoint sets in \(\mathcal{S}\) whose union is contained in \(X\), so the definition 9.8 of \(|\nu|\) gives

\begin{equation*} |\nu(E)| + |\nu(X \setminus E)| \le |\nu|(X) = \nu(X) = \nu(E) + \nu(X \setminus E), \end{equation*}

where the last equality uses the additivity of \(\nu\). Because the left side above is a real number and \(\nu(X)\) is real, we may take real parts on the right, obtaining

\begin{equation*} |\nu(E)| + |\nu(X \setminus E)| \le \operatorname{Re} \nu(E) + \operatorname{Re} \nu(X \setminus E). \end{equation*}

On the other hand, \(\operatorname{Re} z \le |z|\) for every \(z \in \mathbb{C}\), so

\begin{equation*} \operatorname{Re} \nu(E) + \operatorname{Re} \nu(X \setminus E) \le |\nu(E)| + |\nu(X \setminus E)|. \end{equation*}

The two displayed inequalities force equality throughout. Because \(\operatorname{Re} \nu(E) \le |\nu(E)|\) and \(\operatorname{Re} \nu(X \setminus E) \le |\nu(X \setminus E)|\), equality of the sums forces equality termwise:

\begin{equation*} \operatorname{Re} \nu(E) = |\nu(E)|. \end{equation*}

For a complex number \(z\), the equation \(\operatorname{Re} z = |z|\) holds if and only if \(z \in [0, \infty)\) (indeed \(\operatorname{Re} z = |z|\) implies \((\operatorname{Re} z)^2 = (\operatorname{Re} z)^2 + (\operatorname{Im} z)^2\), so \(\operatorname{Im} z = 0\), and then \(z = \operatorname{Re} z = |z| \ge 0\)). Hence

\begin{equation*} \nu(E) \in [0, \infty). \end{equation*}

Because \(E \in \mathcal{S}\) was arbitrary, \(\nu\) takes values in \([0, \infty)\). Also \(\nu(\varnothing) = 0\) by 9.3(a), and \(\nu\) is countably additive because it is a complex measure. Therefore \(\nu\) is a (positive) measure (indeed a finite one), completing the proof.

Problem (9A.4)

Suppose \(\nu\) is a complex measure on a measurable space \((X, \mathcal{S})\). Prove that if \(E \in \mathcal{S}\) then

\begin{equation*} |\nu|(E) = \sup\Bigl\{ \sum_{k=1}^{\infty} |\nu(E_k)| : E_1, E_2, \ldots \text{ is a disjoint sequence in } \mathcal{S} \text{ such that } E = \bigcup_{k=1}^{\infty} E_k \Bigr\}. \end{equation*}

Solution

Fix \(E \in \mathcal{S}\) and let \(s\) denote the supremum on the right side of the displayed equation, that is,

\begin{equation*} s = \sup\Bigl\{ \sum_{k=1}^{\infty} |\nu(E_k)| : E_1, E_2, \ldots \text{ is a disjoint sequence in } \mathcal{S} \text{ with } \bigcup_{k=1}^{\infty} E_k = E \Bigr\}. \end{equation*}

The set over which this supremum is taken is nonempty: the sequence \(E, \varnothing, \varnothing, \ldots\) is a disjoint sequence in \(\mathcal{S}\) whose union is \(E\). Also each sum \(\sum_{k=1}^{\infty} |\nu(E_k)|\) appearing above is a well-defined element of \([0, \infty)\), by 9.3(b). Thus \(s\) is a well-defined element of \([0, \infty]\).

We prove the two inequalities \(s \le |\nu|(E)\) and \(|\nu|(E) \le s\).

Proof that \(s \le |\nu|(E)\). Suppose \(E_1, E_2, \ldots\) is a disjoint sequence in \(\mathcal{S}\) with \(\bigcup_{k=1}^{\infty} E_k = E\). Let \(n \in \mathbb{Z}^{+}\). Then \(E_1, \ldots, E_n\) are disjoint sets in \(\mathcal{S}\) with

\begin{equation*} E_1 \cup \cdots \cup E_n \subseteq \bigcup_{k=1}^{\infty} E_k = E, \end{equation*}

so the definition 9.8 of the total variation measure gives

\begin{equation*} \sum_{k=1}^{n} |\nu(E_k)| \le |\nu|(E). \end{equation*}

Because this holds for every \(n \in \mathbb{Z}^{+}\), letting \(n \to \infty\) gives

\begin{equation*} \sum_{k=1}^{\infty} |\nu(E_k)| \le |\nu|(E). \end{equation*}

Taking the supremum over all such disjoint sequences shows that \(s \le |\nu|(E)\).

Proof that \(|\nu|(E) \le s\). Suppose \(n \in \mathbb{Z}^{+}\) and \(E_1, \ldots, E_n\) are disjoint sets in \(\mathcal{S}\) such that \(E_1 \cup \cdots \cup E_n \subseteq E\). Define a sequence \(F_1, F_2, \ldots\) of sets in \(\mathcal{S}\) by

\begin{equation*} F_k = E_k \ \text{ for } k \in \{1, \ldots, n\}, \qquad F_{n+1} = E \setminus (E_1 \cup \cdots \cup E_n), \qquad F_k = \varnothing \ \text{ for } k > n+1. \end{equation*}

Each \(F_k\) is in \(\mathcal{S}\) because \(\mathcal{S}\) is a \(\sigma\)-algebra. The sets \(F_1, F_2, \ldots\) are pairwise disjoint: the sets \(E_1, \ldots, E_n\) are pairwise disjoint by hypothesis, the set \(F_{n+1}\) is disjoint from each of \(E_1, \ldots, E_n\) by construction, and the remaining sets are empty. Furthermore,

\begin{equation*} \bigcup_{k=1}^{\infty} F_k = (E_1 \cup \cdots \cup E_n) \cup \bigl( E \setminus (E_1 \cup \cdots \cup E_n) \bigr) = E, \end{equation*}

where the last equality uses \(E_1 \cup \cdots \cup E_n \subseteq E\).

Thus \(F_1, F_2, \ldots\) is one of the disjoint sequences appearing in the supremum defining \(s\), and therefore

\begin{equation*} \sum_{k=1}^{n} |\nu(E_k)| \le \sum_{k=1}^{n} |\nu(E_k)| + |\nu(F_{n+1})| = \sum_{k=1}^{\infty} |\nu(F_k)| \le s, \end{equation*}

where the middle equality holds because \(\nu(F_k) = \nu(\varnothing) = 0\) for \(k > n+1\), by 9.3(a).

We have shown that \(\sum_{k=1}^{n} |\nu(E_k)| \le s\) for every \(n \in \mathbb{Z}^{+}\) and all disjoint \(E_1, \ldots, E_n \in \mathcal{S}\) with \(E_1 \cup \cdots \cup E_n \subseteq E\). Taking the supremum over all such choices and using the definition 9.8 of \(|\nu|(E)\) gives \(|\nu|(E) \le s\).

Combining the two inequalities gives \(|\nu|(E) = s\), as desired.

Problem (9A.5)

Suppose \(\mu\) is a (positive) measure on a measurable space \((X, \mathcal{S})\) and \(h\) is a nonnegative function in \(\mathcal{L}^1(\mu)\). Let \(\nu\) be the (positive) measure on \((X, \mathcal{S})\) defined by \(d\nu = h \, d\mu\). Prove that

\begin{equation*} \int f \, d\nu = \int f h \, d\mu \end{equation*}

for all \(\mathcal{S}\)-measurable functions \(f : X \to [0, \infty]\).

Solution

Recall that \(d\nu = h \, d\mu\) means

\begin{equation*} \nu(E) = \int_E h \, d\mu \qquad \text{for every } E \in \mathcal{S}; \end{equation*}

this is a complex (indeed real) measure by 9.4, and because \(h\) is nonnegative we have \(\nu(E) \ge 0\) for every \(E \in \mathcal{S}\), so \(\nu\) is a finite (positive) measure. Because \(h \in \mathcal{L}^1(\mu)\) is nonnegative, \(h\) is an \(\mathcal{S}\)-measurable function from \(X\) to \([0, \infty)\). Thus for \(\mathcal{S}\)-measurable \(f : X \to [0, \infty]\) the product \(f h : X \to [0, \infty]\) is \(\mathcal{S}\)-measurable, using the convention \(0 \cdot \infty = 0\).

We proceed by the standard three-step argument: characteristic functions, then nonnegative simple functions, then general nonnegative measurable functions.

Step 1: characteristic functions. Suppose \(A \in \mathcal{S}\) and \(f = \chi_A\). Applying 3.15 (with \(n = 1\) and \(c_1 = 1\)) with respect to the measure \(\nu\) gives

\begin{equation*} \int \chi_A \, d\nu = \nu(A). \end{equation*}

On the other hand, by the definition of \(\nu\) and by 3.24 (integration on a subset),

\begin{equation*} \nu(A) = \int_A h \, d\mu = \int \chi_A h \, d\mu. \end{equation*}

Combining the two displayed equations gives \(\int \chi_A \, d\nu = \int \chi_A h \, d\mu\), which is the desired conclusion for \(f = \chi_A\).

Step 2: nonnegative simple functions. Suppose \(g : X \to [0, \infty)\) is a simple \(\mathcal{S}\)-measurable function. If \(X = \varnothing\) then every integral below equals \(0\) and there is nothing to prove, so assume \(X \ne \varnothing\). Because \(g\) is simple (2.88), it takes only finitely many values; let \(c_1, \ldots, c_n \in [0, \infty)\) be the distinct values of \(g\) and put \(A_k = g^{-1}(\{c_k\})\). Each \(A_k\) is in \(\mathcal{S}\) because \(g\) is \(\mathcal{S}\)-measurable, the sets \(A_1, \ldots, A_n\) are disjoint with union \(X\), and

\begin{equation*} g = \sum_{k=1}^{n} c_k \chi_{A_k} \end{equation*}

(this is the standard representation discussed just after 2.88). Then 3.15 applied with respect to \(\nu\) gives

\begin{equation*} \int g \, d\nu = \sum_{k=1}^{n} c_k \, \nu(A_k). \end{equation*}

On the other side, \(g h = \sum_{k=1}^{n} c_k \chi_{A_k} h\) is a sum of \(n\) nonnegative \(\mathcal{S}\)-measurable functions, so \(n - 1\) applications of the additivity of integration (3.16), followed by the homogeneity of integration (3.20) applied to each term, give

\begin{equation*} \int g h \, d\mu = \sum_{k=1}^{n} \int c_k \chi_{A_k} h \, d\mu = \sum_{k=1}^{n} c_k \int \chi_{A_k} h \, d\mu = \sum_{k=1}^{n} c_k \, \nu(A_k), \end{equation*}

where the last equality uses Step 1. Comparing the last two displayed equations gives

\begin{equation*} \int g \, d\nu = \int g h \, d\mu \end{equation*}

for every nonnegative simple \(\mathcal{S}\)-measurable function \(g\).

Step 3: general nonnegative measurable functions. Suppose \(f : X \to [0, \infty]\) is \(\mathcal{S}\)-measurable. By 2.89 there exists a sequence \(g_1, g_2, \ldots\) of simple \(\mathcal{S}\)-measurable functions from \(X\) to \(\mathbb{R}\) such that \(|g_k(x)| \le |g_{k+1}(x)| \le |f(x)|\) for all \(k \in \mathbb{Z}^{+}\) and all \(x \in X\) (2.89(b)) and \(\lim_{k \to \infty} g_k(x) = f(x)\) for every \(x \in X\) (2.89(c)).

Set \(f_k = |g_k|\). Each \(f_k\) is a simple \(\mathcal{S}\)-measurable function from \(X\) to \([0, \infty)\), and because \(f\) takes values in \([0, \infty]\) we have \(|f| = f\). Hence

\begin{equation*} 0 \le f_1(x) \le f_2(x) \le \cdots \le f(x) \quad \text{and} \quad \lim_{k \to \infty} f_k(x) = |f(x)| = f(x) \end{equation*}

for every \(x \in X\), where the limit statement uses the continuity of the absolute value on \([-\infty, \infty]\). This avoids having to know that the functions produced by 2.89 are themselves nonnegative.

Because \(h(x) \ge 0\) for every \(x \in X\), multiplying by \(h(x)\) preserves these inequalities, so

\begin{equation*} 0 \le f_1(x) h(x) \le f_2(x) h(x) \le \cdots \end{equation*}

for every \(x \in X\). Moreover \(\lim_{k \to \infty} f_k(x) h(x) = f(x) h(x)\) for every \(x \in X\): if \(h(x) = 0\) then both sides are \(0\) (using the convention \(0 \cdot \infty = 0\)), and if \(h(x) > 0\) then \(f_k(x) h(x) \to f(x) h(x)\) in \([0, \infty]\) because \(f_k(x) \to f(x)\).

Now apply the Monotone Convergence Theorem (3.11) twice. With respect to the measure \(\nu\) it gives

\begin{equation*} \lim_{k \to \infty} \int f_k \, d\nu = \int f \, d\nu, \end{equation*}

and with respect to the measure \(\mu\), applied to the increasing sequence \(f_1 h, f_2 h, \ldots\) of nonnegative \(\mathcal{S}\)-measurable functions converging pointwise to \(f h\), it gives

\begin{equation*} \lim_{k \to \infty} \int f_k h \, d\mu = \int f h \, d\mu. \end{equation*}

By Step 2, \(\int f_k \, d\nu = \int f_k h \, d\mu\) for every \(k \in \mathbb{Z}^{+}\). Hence the two limits above are limits of the same sequence, and therefore

\begin{equation*} \int f \, d\nu = \int f h \, d\mu, \end{equation*}

completing the proof.

Problem (9A.6)

Suppose \((X, \mathcal{S}, \mu)\) is a (positive) measure space. Prove that

\begin{equation*} \{ h \, d\mu : h \in \mathcal{L}^1(\mu) \} \end{equation*}

is a closed subspace of \(M_{\mathbf{F}}(\mathcal{S})\).

Solution

Throughout, \(h \, d\mu\) denotes the function on \(\mathcal{S}\) defined by

\begin{equation*} (h \, d\mu)(E) = \int_E h \, d\mu, \end{equation*}

which is an element of \(M_{\mathbf{F}}(\mathcal{S})\) by 9.4. Set

\begin{equation*} M = \{ h \, d\mu : h \in \mathcal{L}^1(\mu) \} \subseteq M_{\mathbf{F}}(\mathcal{S}). \end{equation*}

Step 1: \(M\) is a subspace of \(M_{\mathbf{F}}(\mathcal{S})\).

The zero measure is in \(M\), because \(0 \, d\mu\) is the measure that assigns \(0\) to every set in \(\mathcal{S}\) and the zero function is in \(\mathcal{L}^1(\mu)\). Thus \(M \ne \varnothing\).

Suppose \(h, g \in \mathcal{L}^1(\mu)\) and \(\alpha \in \mathbf{F}\). Then \(h + \alpha g \in \mathcal{L}^1(\mu)\), because \(\mathcal{L}^1(\mu)\) is a vector space (7.5 with \(p = 1\)). If \(E \in \mathcal{S}\), then the linearity of integration (3.20 and 3.21) gives

\begin{equation*} \bigl( h \, d\mu + \alpha (g \, d\mu) \bigr)(E) = \int_E h \, d\mu + \alpha \int_E g \, d\mu = \int_E (h + \alpha g) \, d\mu = \bigl( (h + \alpha g) \, d\mu \bigr)(E), \end{equation*}

where the definition 9.13 of addition and scalar multiplication of measures was used on the left. Hence

\begin{equation*} h \, d\mu + \alpha (g \, d\mu) = (h + \alpha g) \, d\mu \in M. \end{equation*}

Thus \(M\) is closed under addition and scalar multiplication, so \(M\) is a subspace of \(M_{\mathbf{F}}(\mathcal{S})\).

Step 2: the key isometry.

If \(h \in \mathcal{L}^1(\mu)\), then 9.10 applied with \(E = X\) gives

\begin{equation*} \| h \, d\mu \| = |h \, d\mu|(X) = \int_X |h| \, d\mu = \|h\|_1; \end{equation*}

this is the second bullet point of 9.16. Because \(M\) is a subspace and because the assignment \(h \mapsto h \, d\mu\) is linear (Step 1), we get the crucial identity

9A.6.a: \(\| h \, d\mu - g \, d\mu \| = \|(h - g) \, d\mu\| = \|h - g\|_1\) for all \(h, g \in \mathcal{L}^1(\mu)\).

Step 3: \(M\) is closed.

Suppose \(\nu_1, \nu_2, \ldots\) is a sequence in \(M\) and \(\nu \in M_{\mathbf{F}}(\mathcal{S})\) with \(\lim_{n \to \infty} \|\nu_n - \nu\| = 0\). We must show that \(\nu \in M\).

For each \(n \in \mathbb{Z}^{+}\), choose \(h_n \in \mathcal{L}^1(\mu)\) with \(\nu_n = h_n \, d\mu\). Because the sequence \(\nu_1, \nu_2, \ldots\) converges in the normed space \(M_{\mathbf{F}}(\mathcal{S})\), it is a Cauchy sequence there. Hence by 9A.6.a,

\begin{equation*} \|h_j - h_k\|_1 = \|\nu_j - \nu_k\| \to 0 \qquad \text{as } j, k \to \infty; \end{equation*}

that is, given \(\varepsilon > 0\) there exists \(m \in \mathbb{Z}^{+}\) such that \(\|h_j - h_k\|_1 < \varepsilon\) for all \(j, k \ge m\).

Consider now the corresponding elements \(\tilde{h}_1, \tilde{h}_2, \ldots\) of the Lebesgue space \(L^1(\mu)\); here 7.15 and 7.16 are in force, so \(\tilde{h}\) denotes the set of functions in \(\mathcal{L}^1(\mu)\) that agree with \(h\) almost everywhere, and by 7.17 the norm on \(L^1(\mu)\) satisfies

\begin{equation*} \|\tilde{u} - \tilde{v}\|_1 = \|\widetilde{u - v}\|_1 = \|u - v\|_1 \qquad \text{for all } u, v \in \mathcal{L}^1(\mu). \end{equation*}

Hence the displayed Cauchy condition on \(\|h_j - h_k\|_1\) says exactly that \(\tilde{h}_1, \tilde{h}_2, \ldots\) is a Cauchy sequence in \(L^1(\mu)\). Because \(L^1(\mu)\) is a Banach space (7.24), this sequence converges to some element of \(L^1(\mu)\), which by the definition 7.16 has the form \(\tilde{h}\) for some \(h \in \mathcal{L}^1(\mu)\). Using the displayed norm identity again,

\begin{equation*} \lim_{n \to \infty} \|h_n - h\|_1 = \lim_{n \to \infty} \|\tilde{h}_n - \tilde{h}\|_1 = 0. \end{equation*}

Now 9A.6.a gives

\begin{equation*} \|\nu_n - h \, d\mu\| = \|h_n \, d\mu - h \, d\mu\| = \|h_n - h\|_1 \to 0, \end{equation*}

so \(\nu_n \to h \, d\mu\) in \(M_{\mathbf{F}}(\mathcal{S})\). But also \(\nu_n \to \nu\). Limits in a normed vector space are unique, because

\begin{equation*} \|\nu - h \, d\mu\| \le \|\nu - \nu_n\| + \|\nu_n - h \, d\mu\| \to 0 \end{equation*}

forces \(\|\nu - h \, d\mu\| = 0\) and hence \(\nu = h \, d\mu\) (recall from the remarks following 9.17 that the total variation norm really is a norm on \(M_{\mathbf{F}}(\mathcal{S})\)).

Therefore \(\nu = h \, d\mu \in M\), which shows that \(M\) contains all limits of convergent sequences of its elements. Hence \(M\) is a closed subspace of \(M_{\mathbf{F}}(\mathcal{S})\).

Remark. The argument shows more: by 9A.6.a the map \(\tilde{h} \mapsto h \, d\mu\) is a well-defined linear isometry of the Banach space \(L^1(\mu)\) onto \(M\) (well defined because \(\|h - g\|_1 = 0\) implies \(h \, d\mu = g \, d\mu\), by 9A.6.a and the fact that the total variation norm is a norm). Thus \(M\) is the isometric image of a complete space, hence complete, hence closed in \(M_{\mathbf{F}}(\mathcal{S})\). In the language of Chapter 9, the closed subspace \(M\) consists exactly of the measures in \(M_{\mathbf{F}}(\mathcal{S})\) that are absolutely continuous with respect to \(\mu\) when \(\mu\) is \(\sigma\)-finite; that identification is the Radon-Nikodym Theorem and is not needed here.

Problem (9A.7)

(a) Suppose \(\mathcal{B}\) is the collection of Borel subsets of \(\mathbb{R}\). Show that the Banach space \(M_{\mathbf{F}}(\mathcal{B})\) is not separable.

(b) Give an example of a measurable space \((X, \mathcal{S})\) such that the Banach space \(M_{\mathbf{F}}(\mathcal{S})\) is infinite-dimensional and separable.

Solution

Recall that a metric space is called separable if it has a countable dense subset.

Part (a).

For \(t \in \mathbb{R}\), let \(\delta_t : \mathcal{B} \to \mathbf{F}\) be the Dirac measure at \(t\), defined by

\begin{equation*} \delta_t(E) = \begin{cases} 1 & \text{if } t \in E, \\ 0 & \text{if } t \notin E. \end{cases} \end{equation*}

Then \(\delta_t \in M_{\mathbf{F}}(\mathcal{B})\). To see this, suppose \(E_1, E_2, \ldots\) is a disjoint sequence of Borel sets with union \(E\). If \(t \in E\), then \(t\) belongs to exactly one of the sets \(E_k\) (exactly one because the sets are disjoint), so \(\sum_{k=1}^{\infty} \delta_t(E_k) = 1 = \delta_t(E)\); if \(t \notin E\), then \(t\) belongs to none of the sets \(E_k\), so \(\sum_{k=1}^{\infty} \delta_t(E_k) = 0 = \delta_t(E)\). Thus \(\delta_t\) is countably additive and takes values in \(\{0, 1\} \subseteq \mathbf{F}\), so \(\delta_t\) is a real measure (and hence also a complex measure).

Now suppose \(s, t \in \mathbb{R}\) with \(s \ne t\). The singletons \(\{s\}\) and \(\{t\}\) are disjoint Borel sets, and

\begin{equation*} (\delta_s - \delta_t)(\{s\}) = 1 - 0 = 1, \qquad (\delta_s - \delta_t)(\{t\}) = 0 - 1 = -1. \end{equation*}

Hence the definition 9.8 of the total variation measure (taking \(n = 2\), \(E_1 = \{s\}\), \(E_2 = \{t\}\)) gives

\begin{equation*} \|\delta_s - \delta_t\| = |\delta_s - \delta_t|(\mathbb{R}) \ge |1| + |-1| = 2. \end{equation*}

Suppose now, seeking a contradiction, that \(M_{\mathbf{F}}(\mathcal{B})\) is separable, and let \(D\) be a countable dense subset of \(M_{\mathbf{F}}(\mathcal{B})\). For each \(t \in \mathbb{R}\), density gives some \(\nu_t \in D\) with

\begin{equation*} \|\nu_t - \delta_t\| < 1. \end{equation*}

If \(s, t \in \mathbb{R}\) with \(s \ne t\) and \(\nu_s = \nu_t\), then the triangle inequality would give

\begin{equation*} 2 \le \|\delta_s - \delta_t\| \le \|\delta_s - \nu_s\| + \|\nu_t - \delta_t\| < 1 + 1 = 2, \end{equation*}

a contradiction. Hence \(s \mapsto \nu_s\) is an injective function from \(\mathbb{R}\) into \(D\), which implies that \(D\) is uncountable. This contradicts the countability of \(D\).

Therefore \(M_{\mathbf{F}}(\mathcal{B})\) is not separable.

Part (b).

Take \(X = \mathbb{Z}^{+}\) and let \(\mathcal{S}\) be the \(\sigma\)-algebra of all subsets of \(\mathbb{Z}^{+}\). We show that \(M_{\mathbf{F}}(\mathcal{S})\) is isometrically isomorphic to \(\ell^1\) (see 7.4), and that consequently it is infinite-dimensional and separable.

Define \(T : M_{\mathbf{F}}(\mathcal{S}) \to \ell^1\) by

\begin{equation*} T\nu = \bigl( \nu(\{1\}), \nu(\{2\}), \ldots \bigr). \end{equation*}

\(T\) maps into \(\ell^1\): the sets \(\{1\}, \{2\}, \ldots\) form a disjoint sequence in \(\mathcal{S}\), so \(\sum_{k=1}^{\infty} |\nu(\{k\})| < \infty\) by 9.3(b).

\(T\) is linear: immediate from the definition 9.13 of addition and scalar multiplication of measures.

\(T\) is injective: if \(E \in \mathcal{S}\), then \(E = \bigcup_{k \in E} \{k\}\) is a countable disjoint union of sets in \(\mathcal{S}\) (pad with copies of \(\varnothing\) if \(E\) is finite), so countable additivity gives

\begin{equation*} \nu(E) = \sum_{k \in E} \nu(\{k\}). \end{equation*}

Thus \(T\nu = 0\) implies \(\nu(E) = 0\) for every \(E \in \mathcal{S}\), so \(\nu = 0\).

\(T\) is surjective: suppose \(a = (a_1, a_2, \ldots) \in \ell^1\). Define \(\nu : \mathcal{S} \to \mathbf{F}\) by

\begin{equation*} \nu(E) = \sum_{k \in E} a_k. \end{equation*}

This series converges absolutely (it is a subseries of an absolutely convergent series), so \(\nu(E)\) is a well-defined element of \(\mathbf{F}\), independent of the order of summation. If \(E_1, E_2, \ldots\) is a disjoint sequence in \(\mathcal{S}\) with union \(E\), then the double family \(\{a_k : k \in E_j,\ j \in \mathbb{Z}^{+}\}\) is an enumeration without repetition of \(\{a_k : k \in E\}\), and because \(\sum_{k \in E} |a_k| < \infty\) an absolutely convergent series may be grouped and rearranged arbitrarily without changing its sum. Hence

\begin{equation*} \sum_{j=1}^{\infty} \nu(E_j) = \sum_{j=1}^{\infty} \sum_{k \in E_j} a_k = \sum_{k \in E} a_k = \nu(E). \end{equation*}

Thus \(\nu\) is countably additive, so \(\nu \in M_{\mathbf{F}}(\mathcal{S})\), and clearly \(T\nu = a\).

\(T\) is an isometry: suppose \(\nu \in M_{\mathbf{F}}(\mathcal{S})\) and \(a = T\nu\). If \(n \in \mathbb{Z}^{+}\) and \(E_1, \ldots, E_n\) are disjoint sets in \(\mathcal{S}\), then

\begin{equation*} \sum_{j=1}^{n} |\nu(E_j)| = \sum_{j=1}^{n} \Bigl| \sum_{k \in E_j} a_k \Bigr| \le \sum_{j=1}^{n} \sum_{k \in E_j} |a_k| \le \sum_{k=1}^{\infty} |a_k| = \|a\|_1, \end{equation*}

where the last inequality holds because the sets \(E_1, \ldots, E_n\) are disjoint. Taking the supremum over all such choices gives \(\|\nu\| = |\nu|(\mathbb{Z}^{+}) \le \|a\|_1\). For the reverse inequality, fix \(n \in \mathbb{Z}^{+}\) and take \(E_j = \{j\}\) for \(j \in \{1, \ldots, n\}\) in the definition 9.8 of \(|\nu|(\mathbb{Z}^{+})\); this gives

\begin{equation*} \sum_{j=1}^{n} |a_j| = \sum_{j=1}^{n} |\nu(\{j\})| \le \|\nu\|. \end{equation*}

Letting \(n \to \infty\) gives \(\|a\|_1 \le \|\nu\|\). Hence \(\|T\nu\|_1 = \|\nu\|\).

Thus \(T\) is a linear isometric bijection from \(M_{\mathbf{F}}(\mathcal{S})\) onto \(\ell^1\).

\(M_{\mathbf{F}}(\mathcal{S})\) is infinite-dimensional: for \(k \in \mathbb{Z}^{+}\) let \(\delta_k \in M_{\mathbf{F}}(\mathcal{S})\) be the Dirac measure at \(k\) (countably additive by the argument in part (a)). Then \(T\delta_k\) is the sequence with \(1\) in slot \(k\) and \(0\) elsewhere. These sequences are linearly independent in \(\ell^1\) for \(k = 1, 2, \ldots, n\) and any \(n\), so \(\delta_1, \ldots, \delta_n\) are linearly independent in \(M_{\mathbf{F}}(\mathcal{S})\) for every \(n \in \mathbb{Z}^{+}\). Hence \(M_{\mathbf{F}}(\mathcal{S})\) is infinite-dimensional.

\(M_{\mathbf{F}}(\mathcal{S})\) is separable: let \(\mathbf{Q}\) denote \(\mathbb{Q}\) if \(\mathbf{F} = \mathbb{R}\) and \(\{p + iq : p, q \in \mathbb{Q}\}\) if \(\mathbf{F} = \mathbb{C}\); in either case \(\mathbf{Q}\) is a countable dense subset of \(\mathbf{F}\). Let

\begin{equation*} D = \{ (q_1, \ldots, q_n, 0, 0, \ldots) : n \in \mathbb{Z}^{+} \text{ and } q_1, \ldots, q_n \in \mathbf{Q} \}. \end{equation*}

Then \(D \subseteq \ell^1\) and \(D\) is countable, being a countable union (over \(n\)) of the countable sets \(\mathbf{Q}^n\). To see that \(D\) is dense in \(\ell^1\), suppose \(a \in \ell^1\) and \(\varepsilon > 0\). Because \(\sum_{k=1}^{\infty} |a_k| < \infty\), there exists \(n \in \mathbb{Z}^{+}\) with \(\sum_{k=n+1}^{\infty} |a_k| < \varepsilon/2\). Because \(\mathbf{Q}\) is dense in \(\mathbf{F}\), there exist \(q_1, \ldots, q_n \in \mathbf{Q}\) with \(|a_k - q_k| < \varepsilon/(2n)\) for each \(k \in \{1, \ldots, n\}\). Let \(b = (q_1, \ldots, q_n, 0, 0, \ldots) \in D\). Then

\begin{equation*} \|a - b\|_1 = \sum_{k=1}^{n} |a_k - q_k| + \sum_{k=n+1}^{\infty} |a_k| < \frac{\varepsilon}{2} + \frac{\varepsilon}{2} = \varepsilon. \end{equation*}

Thus \(\ell^1\) is separable. Because \(T^{-1}\) is a surjective isometry from \(\ell^1\) onto \(M_{\mathbf{F}}(\mathcal{S})\), the countable set \(T^{-1}(D)\) is dense in \(M_{\mathbf{F}}(\mathcal{S})\): given \(\nu \in M_{\mathbf{F}}(\mathcal{S})\) and \(\varepsilon > 0\), pick \(b \in D\) with \(\|T\nu - b\|_1 < \varepsilon\), and then \(\|\nu - T^{-1}b\| = \|T\nu - b\|_1 < \varepsilon\).

Therefore \(M_{\mathbf{F}}(\mathcal{S})\) is an infinite-dimensional separable Banach space, as desired.

Remark. Comparing the two parts: what makes \(M_{\mathbf{F}}(\mathcal{B})\) huge is the uncountable supply of atoms available in \(\mathbb{R}\), each contributing a Dirac measure at distance \(2\) from all the others. When \(X\) is countable and \(\mathcal{S}\) is its full \(\sigma\)-algebra, the atoms form a countable set and the space collapses to \(\ell^1\).

Problem (9A.8)

Suppose \(t > 0\) and \(\lambda\) is Lebesgue measure on the \(\sigma\)-algebra of Borel subsets of \([0, t]\). Suppose \(h : [0, t] \to \mathbb{C}\) is the function defined by

\begin{equation*} h(x) = \cos x + i \sin x. \end{equation*}

Let \(\nu\) be the complex measure defined by \(d\nu = h \, d\lambda\).

(a) Show that \(\|\nu\| = t\).

(b) Show that if \(E_1, E_2, \ldots\) is a sequence of disjoint Borel subsets of \([0, t]\), then

\begin{equation*} \sum_{k=1}^{\infty} |\nu(E_k)| < t. \end{equation*}

[This exercise shows that the supremum in the definition of \(|\nu|([0,t])\) is not attained, even if countably many disjoint sets are allowed.]

Solution

Throughout, write \(h(x) = e^{ix}\), so that \(|h(x)| = 1\) for every \(x \in [0, t]\). In particular \(h\) is a bounded Borel measurable function on a set of finite measure, so \(h \in L^1(\lambda)\) and \(\nu = h \, d\lambda\) really is a complex measure on the Borel subsets of \([0,t]\) (by 9.4).

(a) By 9.10, for every Borel set \(E \subseteq [0,t]\) we have

\begin{equation*} |\nu|(E) = \int_E |h| \, d\lambda = \int_E 1 \, d\lambda = \lambda(E). \end{equation*}

Thus \(|\nu|\) is just \(\lambda\), and hence

\begin{equation*} \|\nu\| = |\nu|([0,t]) = \lambda([0,t]) = t . \end{equation*}

(b) The key point is the following strict inequality.

Claim: if \(E \subseteq [0,t]\) is a Borel set with \(\lambda(E) > 0\), then \(|\nu(E)| < \lambda(E)\).

To prove the claim, choose \(\alpha \in \mathbb{C}\) with \(|\alpha| = 1\) such that \(\alpha \, \nu(E) = |\nu(E)|\) (take \(\alpha = \overline{\nu(E)}/|\nu(E)|\) if \(\nu(E) \neq 0\), and \(\alpha = 1\) otherwise). Then

\begin{equation*} |\nu(E)| = \alpha \int_E e^{ix} \, d\lambda(x) = \int_E \alpha e^{ix} \, d\lambda(x). \end{equation*}

The left side is real, so the imaginary part of the right side vanishes and

\begin{equation*} |\nu(E)| = \int_E \operatorname{Re}\!\left( \alpha e^{ix} \right) d\lambda(x). \end{equation*}

Because \(|\alpha e^{ix}| = 1\), we have \(\operatorname{Re}(\alpha e^{ix}) \le 1\) for all \(x\), with equality precisely when \(\alpha e^{ix} = 1\). Writing \(\alpha = e^{i\theta}\), equality holds exactly for those \(x\) with \(x + \theta \in 2\pi \mathbb{Z}\), i.e. for \(x\) in a countable set. Hence the Borel set

\begin{equation*} N = \{ x \in E : \operatorname{Re}(\alpha e^{ix}) = 1 \} \end{equation*}

is countable, hence Borel, and \(\lambda(N) = 0\); therefore \(\lambda(E \setminus N) = \lambda(E) > 0\).

Let \(f(x) = 1 - \operatorname{Re}(\alpha e^{ix})\), a Borel measurable function on \(E\) with \(0 \le f \le 2\), so \(f \in L^1(\lambda)\). By construction \(f > 0\) exactly on \(E \setminus N\), so

\begin{equation*} E \setminus N = \bigcup_{n=1}^{\infty} \Bigl\{ x \in E : f(x) \ge \tfrac{1}{n} \Bigr\} . \end{equation*}

Suppose, seeking a contradiction, that \(\int_E f \, d\lambda = 0\). Markov’s inequality (4.1), applied on the measure space \(E\) with the restricted Borel \(\sigma\)-algebra and the restriction of \(\lambda\), and with \(c = 1/n\), gives

\begin{equation*} \lambda\Bigl( \Bigl\{ x \in E : f(x) \ge \tfrac{1}{n} \Bigr\} \Bigr) \le n \int_E f \, d\lambda = 0 \end{equation*}

for every \(n \in \mathbb{Z}^+\). Countable subadditivity of \(\lambda\) then forces \(\lambda(E \setminus N) = 0\), contradicting \(\lambda(E \setminus N) = \lambda(E) > 0\). Therefore

\begin{equation*} \lambda(E) - |\nu(E)| = \int_E \bigl( 1 - \operatorname{Re}(\alpha e^{ix}) \bigr) d\lambda(x) > 0, \end{equation*}

which proves the claim.

Note also the (nonstrict) bound valid for every Borel \(E \subseteq [0,t]\):

\begin{equation*} |\nu(E)| = \left| \int_E h \, d\lambda \right| \le \int_E |h| \, d\lambda = \lambda(E). \end{equation*}

Now suppose \(E_1, E_2, \ldots\) are disjoint Borel subsets of \([0,t]\). By countable additivity and monotonicity of \(\lambda\),

\begin{equation*} \sum_{k=1}^{\infty} \lambda(E_k) = \lambda\!\left( \bigcup_{k=1}^{\infty} E_k \right) \le \lambda([0,t]) = t . \end{equation*}

If \(\lambda(E_k) = 0\) for every \(k\), then \(\nu(E_k) = 0\) for every \(k\) (by the nonstrict bound above), so \(\sum_{k=1}^{\infty} |\nu(E_k)| = 0 < t\), as desired (here \(t > 0\) is used).

Otherwise there exists \(j\) with \(\lambda(E_j) > 0\). Set \(\varepsilon = \lambda(E_j) - |\nu(E_j)|\), which is strictly positive by the claim. Then

\begin{equation*} \begin{aligned} \sum_{k=1}^{\infty} |\nu(E_k)| &= |\nu(E_j)| + \sum_{k \neq j} |\nu(E_k)| \\ &\le \bigl( \lambda(E_j) - \varepsilon \bigr) + \sum_{k \neq j} \lambda(E_k) \\ &= \sum_{k=1}^{\infty} \lambda(E_k) - \varepsilon \\ &\le t - \varepsilon \\ &< t . \end{aligned} \end{equation*}

(All the sums involved are finite, since \(\sum_k \lambda(E_k) \le t < \infty\), so the rearrangement above is legitimate.) This proves (b).

Finally, the bracketed remark: by (a) the total variation \(|\nu|([0,t])\) equals \(t\), yet by (b) no countable disjoint family \(E_1, E_2, \ldots\) of Borel subsets of \([0,t]\) achieves \(\sum_k |\nu(E_k)| = t\). In particular no finite disjoint family does either, so the supremum in Definition 9.8 is approached but never attained.

Problem (9A.9)

Give an example to show that 9.9 can fail if the hypothesis that \(\nu\) is a real measure is replaced by the hypothesis that \(\nu\) is a complex measure.

Solution

Recall that 9.9 asserts, for a real measure \(\nu\) on \((X, \mathcal{S})\) and \(E \in \mathcal{S}\), that

\begin{equation*} |\nu|(E) = \sup \{ |\nu(A)| + |\nu(B)| : A, B \text{ are disjoint sets in } \mathcal{S} \text{ and } A \cup B \subseteq E \}. \end{equation*}

In other words, for real measures it suffices to take \(n = 2\) in Definition 9.8. We give a complex measure for which \(n = 2\) is not enough.

Let \(X = \{1, 2, 3\}\) and let \(\mathcal{S}\) be the \(\sigma\)-algebra of all subsets of \(X\). Put

\begin{equation*} \omega = e^{2\pi i/3} = -\tfrac{1}{2} + \tfrac{\sqrt{3}}{2} i , \end{equation*}

so that \(\omega^3 = 1\), \(\omega \neq 1\), and hence \(1 + \omega + \omega^2 = 0\); note also \(\omega^2 = \overline{\omega}\) and \(|\omega| = |\omega^2| = 1\). Let \(\mu\) be counting measure on \((X, \mathcal{S})\), let \(h : X \to \mathbb{C}\) be defined by

\begin{equation*} h(1) = 1, \qquad h(2) = \omega, \qquad h(3) = \omega^2 , \end{equation*}

and let \(\nu = h \, d\mu\); explicitly,

\begin{equation*} \nu(E) = \sum_{k \in E} h(k) \qquad \text{for } E \subseteq X . \end{equation*}

Since \(X\) is finite, \(h \in L^1(\mu)\), so \(\nu\) is a complex measure on \((X, \mathcal{S})\) by 9.4. (Countable additivity is clear directly: a disjoint sequence of subsets of \(X\) has only finitely many nonempty terms.)

First we compute \(|\nu|(X)\). By 9.10,

\begin{equation*} |\nu|(X) = \int_X |h| \, d\mu = |h(1)| + |h(2)| + |h(3)| = 1 + 1 + 1 = 3 . \end{equation*}

(Equivalently and directly: the disjoint family \(\{1\}, \{2\}, \{3\}\) gives \(|\nu(\{1\})| + |\nu(\{2\})| + |\nu(\{3\})| = 3\), so \(|\nu|(X) \ge 3\); and for any disjoint \(E_1, \ldots, E_n \subseteq X\) the triangle inequality gives \(\sum_{j} |\nu(E_j)| \le \sum_j \sum_{k \in E_j} |h(k)| \le 3\).)

Now we compute \(|\nu(E)|\) for every \(E \subseteq X\):

\begin{equation*} \begin{aligned} &\nu(\emptyset) = 0, \\ &\nu(\{1\}) = 1, \quad \nu(\{2\}) = \omega, \quad \nu(\{3\}) = \omega^2, \\ &\nu(\{1,2\}) = 1 + \omega = -\omega^2, \quad \nu(\{1,3\}) = 1 + \omega^2 = -\omega, \quad \nu(\{2,3\}) = \omega + \omega^2 = -1, \\ &\nu(X) = 1 + \omega + \omega^2 = 0 , \end{aligned} \end{equation*}

where we used \(1 + \omega + \omega^2 = 0\) repeatedly. Since \(|\omega| = |\omega^2| = 1\), we conclude that

\begin{equation*} |\nu(E)| = \begin{cases} 0 & \text{if } E = \emptyset \text{ or } E = X, \\ 1 & \text{otherwise}. \end{cases} \end{equation*}

Thus \(|\nu(E)| \le 1\) for every \(E \in \mathcal{S}\), and therefore

\begin{equation*} \sup \{ |\nu(A)| + |\nu(B)| : A, B \in \mathcal{S} \text{ disjoint}, \ A \cup B \subseteq X \} \le 2 . \end{equation*}

(The value \(2\) is in fact attained, for example by \(A = \{1\}\), \(B = \{2\}\).)

Hence with \(E = X\) the right-hand side of 9.9 equals \(2\), while \(|\nu|(X) = 3\). So 9.9 fails for this complex measure.

The underlying reason is exactly the one that makes the proof of 9.9 work for real measures and break for complex ones: for a real measure the sets \(E_k\) can be grouped by the sign of \(\nu(E_k)\) into just two sets \(A\) and \(B\) without losing any of \(\sum_k |\nu(E_k)|\), because the numbers being added within each group all have the same sign. For complex measures there is no such two-way grouping: here the three values \(1, \omega, \omega^2\) point in three different directions, and any grouping into two sets forces cancellation.

Problem (9A.10)

Suppose \((X, \mathcal{S})\) is a measurable space with \(\mathcal{S} \neq \{\emptyset, X\}\). Prove that the total variation norm on \(M_{\mathbb{F}}(\mathcal{S})\) does not come from an inner product. In other words, show that there does not exist an inner product \(\langle \cdot, \cdot \rangle\) on \(M_{\mathbb{F}}(\mathcal{S})\) such that \(\|\nu\| = \langle \nu, \nu \rangle^{1/2}\) for all \(\nu \in M_{\mathbb{F}}(\mathcal{S})\), where \(\|\cdot\|\) is the usual total variation norm on \(M_{\mathbb{F}}(\mathcal{S})\).

Solution

The strategy is standard: a norm arising from an inner product must satisfy the parallelogram equality, and we exhibit two measures in \(M_{\mathbb{F}}(\mathcal{S})\) violating it.

Step 1: the parallelogram equality. Suppose \(\langle \cdot, \cdot \rangle\) is an inner product on a vector space \(V\) over \(\mathbb{F}\) and \(\|u\| = \langle u, u \rangle^{1/2}\) for all \(u \in V\). Then for all \(u, v \in V\), additivity of the inner product in each slot gives

\begin{equation*} \begin{aligned} \|u + v\|^2 + \|u - v\|^2 &= \langle u + v, u + v \rangle + \langle u - v, u - v \rangle \\ &= \bigl( \langle u,u \rangle + \langle u,v \rangle + \langle v,u \rangle + \langle v,v \rangle \bigr) + \bigl( \langle u,u \rangle - \langle u,v \rangle - \langle v,u \rangle + \langle v,v \rangle \bigr) \\ &= 2\|u\|^2 + 2\|v\|^2 . \end{aligned} \end{equation*}

So it suffices to produce \(\nu, \eta \in M_{\mathbb{F}}(\mathcal{S})\) with \(\|\nu + \eta\|^2 + \|\nu - \eta\|^2 \neq 2\|\nu\|^2 + 2\|\eta\|^2\).

Step 2: choosing the two measures. Because \(\mathcal{S} \neq \{\emptyset, X\}\), there exists \(E \in \mathcal{S}\) with \(E \neq \emptyset\) and \(E \neq X\). (Note \(X \neq \emptyset\): if \(X = \emptyset\) then the only \(\sigma\)-algebra on \(X\) is \(\{\emptyset\} = \{\emptyset, X\}\).) Pick \(b \in E\) and \(c \in X \setminus E\); in particular \(b \neq c\).

For \(a \in X\) let \(\delta_a\) denote the Dirac measure at \(a\), defined on \(\mathcal{S}\) by \(\delta_a(A) = 1\) if \(a \in A\) and \(\delta_a(A) = 0\) if \(a \notin A\). This is a finite (positive) measure: \(\delta_a(\emptyset) = 0\), and if \(A_1, A_2, \ldots\) are disjoint sets in \(\mathcal{S}\) then \(a\) belongs to at most one of them, so \(\delta_a(\bigcup_k A_k) = \sum_k \delta_a(A_k)\) (both sides equal \(1\) if \(a\) lies in one of the \(A_k\), and \(0\) otherwise). In particular \(\delta_b, \delta_c\) are real measures, so they belong to \(M_{\mathbb{F}}(\mathcal{S})\) whether \(\mathbb{F} = \mathbb{R}\) or \(\mathbb{F} = \mathbb{C}\).

Set \(\nu = \delta_b\) and \(\eta = \delta_c\).

Step 3: computing the four norms. If \(\mu\) is a finite (positive) measure on \((X,\mathcal{S})\), then \(\|\mu\| = \mu(X)\) (see 9.16): indeed, for disjoint \(E_1, \ldots, E_n \in \mathcal{S}\) with \(E_1 \cup \cdots \cup E_n \subseteq X\) we have \(\sum_{k=1}^n |\mu(E_k)| = \sum_{k=1}^n \mu(E_k) = \mu(E_1 \cup \cdots \cup E_n) \le \mu(X)\), while the single set \(E_1 = X\) gives the value \(\mu(X)\).

Since \(\delta_b\), \(\delta_c\), and \(\delta_b + \delta_c\) are all finite positive measures,

\begin{equation*} \|\nu\| = \delta_b(X) = 1, \qquad \|\eta\| = \delta_c(X) = 1, \qquad \|\nu + \eta\| = (\delta_b + \delta_c)(X) = 2 . \end{equation*}

For \(\nu - \eta = \delta_b - \delta_c\) we bound in both directions. For the upper bound, argue directly from Definition 9.8 rather than invoking the triangle inequality for a norm not yet established: if \(E_1, \ldots, E_n \in \mathcal{S}\) are disjoint, then

\begin{equation*} \sum_{k=1}^{n} \bigl| (\delta_b - \delta_c)(E_k) \bigr| \le \sum_{k=1}^{n} \delta_b(E_k) + \sum_{k=1}^{n} \delta_c(E_k) \le \delta_b(X) + \delta_c(X) = 2 , \end{equation*}

using finite additivity and monotonicity of the positive measures \(\delta_b\) and \(\delta_c\). Hence \(\|\delta_b - \delta_c\| \le 2\). For the lower bound, \(E\) and \(X \setminus E\) are disjoint sets in \(\mathcal{S}\) with \(b \in E\) and \(c \in X \setminus E\), so

\begin{equation*} (\delta_b - \delta_c)(E) = 1 - 0 = 1, \qquad (\delta_b - \delta_c)(X \setminus E) = 0 - 1 = -1 , \end{equation*}

and hence by Definition 9.8 (taking \(n = 2\), \(E_1 = E\), \(E_2 = X \setminus E\))

\begin{equation*} \|\delta_b - \delta_c\| = |\delta_b - \delta_c|(X) \ge |1| + |-1| = 2 . \end{equation*}

Therefore \(\|\nu - \eta\| = 2\).

Step 4: the contradiction. With these values,

\begin{equation*} \|\nu + \eta\|^2 + \|\nu - \eta\|^2 = 2^2 + 2^2 = 8 , \end{equation*}

while

\begin{equation*} 2\|\nu\|^2 + 2\|\eta\|^2 = 2 \cdot 1 + 2 \cdot 1 = 4 . \end{equation*}

Since \(8 \neq 4\), the parallelogram equality fails for the total variation norm on \(M_{\mathbb{F}}(\mathcal{S})\). By Step 1, no inner product \(\langle \cdot, \cdot \rangle\) on \(M_{\mathbb{F}}(\mathcal{S})\) can satisfy \(\|\nu\| = \langle \nu, \nu \rangle^{1/2}\) for all \(\nu \in M_{\mathbb{F}}(\mathcal{S})\).

Remark: the hypothesis \(\mathcal{S} \neq \{\emptyset, X\}\) cannot be dropped. If \(\mathcal{S} = \{\emptyset, X\}\) with \(X \neq \emptyset\), then a measure \(\nu \in M_{\mathbb{F}}(\mathcal{S})\) is determined by the single scalar \(\nu(X)\) and \(\|\nu\| = |\nu(X)|\), so \(M_{\mathbb{F}}(\mathcal{S})\) is isometrically \(\mathbb{F}\) with its absolute value, whose norm does come from an inner product.

Problem (9A.11)

For \((X, \mathcal{S})\) a measurable space and \(b \in X\), define a finite (positive) measure \(\delta_b\) on \((X, \mathcal{S})\) by

\begin{equation*} \delta_b(E) = \begin{cases} 1 & \text{if } b \in E, \\ 0 & \text{if } b \notin E \end{cases} \end{equation*}

for \(E \in \mathcal{S}\).

(a) Show that if \(b, c \in X\), then \(\|\delta_b + \delta_c\| = 2\).

(b) Give an example of a measurable space \((X, \mathcal{S})\) and \(b, c \in X\) with \(b \neq c\) such that \(\|\delta_b - \delta_c\| \neq 2\).

Solution

First note that \(\delta_b\) really is a finite (positive) measure: \(\delta_b(\emptyset) = 0\), and if \(A_1, A_2, \ldots\) are disjoint sets in \(\mathcal{S}\), then \(b\) lies in at most one of them, so both \(\delta_b(\bigcup_k A_k)\) and \(\sum_k \delta_b(A_k)\) equal \(1\) if \(b\) lies in some \(A_k\) and equal \(0\) otherwise. Also \(\delta_b(X) = 1 < \infty\).

We record the fact used in both parts: if \(\mu\) is a finite (positive) measure on \((X, \mathcal{S})\), then

\begin{equation*} \|\mu\| = |\mu|(X) = \mu(X) . \end{equation*}

Indeed, if \(E_1, \ldots, E_n \in \mathcal{S}\) are disjoint with \(E_1 \cup \cdots \cup E_n \subseteq X\), then, using that \(\mu\) is nonnegative and finitely additive and monotone,

\begin{equation*} \sum_{k=1}^{n} |\mu(E_k)| = \sum_{k=1}^{n} \mu(E_k) = \mu(E_1 \cup \cdots \cup E_n) \le \mu(X) , \end{equation*}

so the supremum in Definition 9.8 is at most \(\mu(X)\); taking \(n = 1\) and \(E_1 = X\) shows the supremum is at least \(\mu(X)\). (This is the first bullet point of 9.16.)

(a) Let \(b, c \in X\). The sum \(\delta_b + \delta_c\) is again a finite (positive) measure on \((X, \mathcal{S})\): it is nonnegative, it is countably additive because each summand is, and \((\delta_b + \delta_c)(X) = 1 + 1 = 2 < \infty\). Hence by the fact recorded above,

\begin{equation*} \|\delta_b + \delta_c\| = (\delta_b + \delta_c)(X) = \delta_b(X) + \delta_c(X) = 1 + 1 = 2 . \end{equation*}

Note that this computation is valid whether or not \(b = c\), and regardless of whether \(\mathcal{S}\) distinguishes \(b\) from \(c\): the point is that positive measures cannot cancel, so the total variation norm of a positive measure is just its total mass.

(b) The corresponding statement for the difference is false in general, because \(\delta_b - \delta_c\) can cancel completely when \(\mathcal{S}\) is too coarse to separate \(b\) from \(c\).

Take \(X = \{1, 2\}\) and let \(\mathcal{S} = \{\emptyset, X\}\) be the trivial \(\sigma\)-algebra on \(X\), and set \(b = 1\), \(c = 2\), so \(b \neq c\). Then

\begin{equation*} \delta_b(\emptyset) = 0 = \delta_c(\emptyset), \qquad \delta_b(X) = 1 = \delta_c(X) , \end{equation*}

since \(b \in X\) and \(c \in X\). Thus \(\delta_b(E) = \delta_c(E)\) for every \(E \in \mathcal{S}\), i.e. \(\delta_b - \delta_c\) is the zero measure. Consequently

\begin{equation*} \|\delta_b - \delta_c\| = |0|(X) = 0 \neq 2 . \end{equation*}

Remark: this is the only way the equality \(\|\delta_b - \delta_c\| = 2\) can fail. Suppose \(b, c \in X\). If some \(E \in \mathcal{S}\) contains exactly one of \(b, c\) (say \(b \in E\) and \(c \notin E\)), then \(E\) and \(X \setminus E\) are disjoint sets in \(\mathcal{S}\) with \((\delta_b - \delta_c)(E) = 1\) and \((\delta_b - \delta_c)(X \setminus E) = -1\), so Definition 9.8 (with \(n = 2\)) gives \(\|\delta_b - \delta_c\| \ge |1| + |-1| = 2\). In the other direction, for any disjoint \(E_1, \ldots, E_n \in \mathcal{S}\) we have \(\sum_{k=1}^{n} |(\delta_b - \delta_c)(E_k)| \le \sum_{k=1}^{n} \delta_b(E_k) + \sum_{k=1}^{n} \delta_c(E_k) \le \delta_b(X) + \delta_c(X) = 2\), so \(\|\delta_b - \delta_c\| \le 2\). Hence \(\|\delta_b - \delta_c\| = 2\). If instead no set in \(\mathcal{S}\) separates \(b\) from \(c\), then \(\delta_b(E) = \delta_c(E)\) for all \(E \in \mathcal{S}\) and \(\|\delta_b - \delta_c\| = 0\). So the total variation distance between two Dirac measures is \(2\) when \(\mathcal{S}\) separates the two points and \(0\) when it does not.

Exercises 9B

Problem (9B.1)

Suppose \(\nu\) is a real measure on a measurable space \((X, \mathcal{S})\). Prove that the Hahn decomposition of \(\nu\) is almost unique, in the sense that if \(A, B\) and \(A^{\prime}, B^{\prime}\) are pairs satisfying the Hahn Decomposition Theorem (9.23), then

\begin{equation*} |\nu|(A \setminus A^{\prime}) = |\nu|(A^{\prime} \setminus A) = |\nu|(B \setminus B^{\prime}) = |\nu|(B^{\prime} \setminus B) = 0. \end{equation*}

Solution

First we observe that the four sets appearing in the display are really only two sets. Because \(A \cup B = X\) and \(A \cap B = \varnothing\), we have \(B = X \setminus A\); similarly \(B^{\prime} = X \setminus A^{\prime}\). Hence

\begin{equation*} A \setminus A^{\prime} = A \cap (X \setminus A^{\prime}) = A \cap B^{\prime} = B^{\prime} \setminus B \end{equation*}

and

\begin{equation*} A^{\prime} \setminus A = A^{\prime} \cap (X \setminus A) = A^{\prime} \cap B = B \setminus B^{\prime}. \end{equation*}

Thus \(A \setminus A^{\prime}\) and \(B^{\prime} \setminus B\) are the same set, and \(A^{\prime} \setminus A\) and \(B \setminus B^{\prime}\) are the same set. So it suffices to prove that \(|\nu|(A \setminus A^{\prime}) = 0\) and \(|\nu|(A^{\prime} \setminus A) = 0\).

Suppose \(E \in \mathcal{S}\) and \(E \subseteq A \setminus A^{\prime}\). Then \(E \subseteq A\), so 9.23(b) applied to the Hahn decomposition \(A, B\) gives \(\nu(E) \geq 0\). On the other hand \(E \subseteq B^{\prime}\) (by the first display above), so 9.23(c) applied to the Hahn decomposition \(A^{\prime}, B^{\prime}\) gives \(\nu(E) \leq 0\). Hence

\begin{equation*} \nu(E) = 0 \quad \text{for every } E \in \mathcal{S} \text{ with } E \subseteq A \setminus A^{\prime}. \end{equation*}

Now we compute \(|\nu|(A \setminus A^{\prime})\) directly from the definition 9.8 of the total variation measure. Suppose \(n \in \mathbf{Z}^+\) and \(E_1, \dots, E_n\) are disjoint sets in \(\mathcal{S}\) with \(E_1 \cup \cdots \cup E_n \subseteq A \setminus A^{\prime}\). Then each \(E_k\) is a measurable subset of \(A \setminus A^{\prime}\), so \(\nu(E_k) = 0\) by the previous paragraph, and therefore

\begin{equation*} |\nu(E_1)| + \cdots + |\nu(E_n)| = 0. \end{equation*}

Taking the supremum over all such finite disjoint collections gives \(|\nu|(A \setminus A^{\prime}) = 0\).

The same argument with the roles of the two decompositions interchanged handles the other set: if \(E \in \mathcal{S}\) and \(E \subseteq A^{\prime} \setminus A\), then \(E \subseteq A^{\prime}\) gives \(\nu(E) \geq 0\) by 9.23(b) for \(A^{\prime}, B^{\prime}\), while \(E \subseteq B\) gives \(\nu(E) \leq 0\) by 9.23(c) for \(A, B\); hence \(\nu(E) = 0\), and exactly as above \(|\nu|(A^{\prime} \setminus A) = 0\).

Combining these two conclusions with the set identities from the first paragraph, we get

\begin{equation*} |\nu|(A \setminus A^{\prime}) = |\nu|(A^{\prime} \setminus A) = |\nu|(B \setminus B^{\prime}) = |\nu|(B^{\prime} \setminus B) = 0, \end{equation*}

as desired. In other words, any two Hahn decompositions of \(\nu\) differ only by sets that are null for the total variation measure \(|\nu|\), which is the precise sense in which the Hahn decomposition is almost unique.

Problem (9B.2)

Suppose \(\mu\) is a (positive) measure and \(g, h \in L^1(\mu)\). Prove that \(g \, d\mu \perp h \, d\mu\) if and only if \(g(x)h(x) = 0\) for almost every \(x \in X\).

Solution

Write \(\nu_g\) and \(\nu_h\) for the complex measures on \((X, \mathcal{S})\) defined by \(d\nu_g = g \, d\mu\) and \(d\nu_h = h \, d\mu\); thus \(\nu_g(E) = \int_E g \, d\mu\) and \(\nu_h(E) = \int_E h \, d\mu\) for \(E \in \mathcal{S}\). Throughout, almost every means with respect to \(\mu\).

We first record a standard fact that will be used twice.

Claim. If \(F \in \mathcal{S}\) and \(\int_F |g| \, d\mu = 0\), then \(g(x) = 0\) for almost every \(x \in F\).

To see this, for \(n \in \mathbf{Z}^+\) let \(F_n = \{x \in F : |g(x)| > \tfrac1n\}\), which is in \(\mathcal{S}\) because \(g\) is \(\mathcal{S}\)-measurable. Then

\begin{equation*} \tfrac1n \mu(F_n) \leq \int_{F_n} |g| \, d\mu \leq \int_F |g| \, d\mu = 0, \end{equation*}

so \(\mu(F_n) = 0\) for every \(n\). Because \(\{x \in F : g(x) \neq 0\} = \bigcup_{n=1}^\infty F_n\), countable subadditivity gives \(\mu(\{x \in F : g(x) \neq 0\}) = 0\), proving the claim.

First suppose \(\nu_g \perp \nu_h\). By 9.28 there exist \(A, B \in \mathcal{S}\) with \(A \cup B = X\), \(A \cap B = \varnothing\), and

\begin{equation*} \nu_g(E) = \nu_g(E \cap A) \quad \text{and} \quad \nu_h(E) = \nu_h(E \cap B) \quad \text{for all } E \in \mathcal{S}. \end{equation*}

Suppose \(F \in \mathcal{S}\) and \(F \subseteq B\). Then \(F \cap A = \varnothing\), so the first equation above gives

\begin{equation*} \nu_g(F) = \nu_g(F \cap A) = \nu_g(\varnothing) = 0. \end{equation*}

Thus \(\nu_g(F) = 0\) for every \(F \in \mathcal{S}\) with \(F \subseteq B\); by the third bullet point following 9.8, this says exactly that \(|\nu_g|(B) = 0\). Now 9.10 tells us that \(|\nu_g|(B) = \int_B |g| \, d\mu\), so \(\int_B |g| \, d\mu = 0\), and the claim gives \(g = 0\) almost everywhere on \(B\). Symmetrically, every \(F \in \mathcal{S}\) with \(F \subseteq A\) satisfies \(\nu_h(F) = \nu_h(F \cap B) = 0\), so \(|\nu_h|(A) = 0\), so \(\int_A |h| \, d\mu = 0\) by 9.10, so \(h = 0\) almost everywhere on \(A\).

Let \(N_1 = \{x \in B : g(x) \neq 0\}\) and \(N_2 = \{x \in A : h(x) \neq 0\}\); we have just shown \(\mu(N_1) = \mu(N_2) = 0\), so \(\mu(N_1 \cup N_2) = 0\). If \(x \in X \setminus (N_1 \cup N_2)\), then either \(x \in A\), in which case \(h(x) = 0\), or \(x \in B\), in which case \(g(x) = 0\); either way \(g(x)h(x) = 0\). Hence \(g(x)h(x) = 0\) for almost every \(x \in X\).

Conversely, suppose \(g(x)h(x) = 0\) for almost every \(x \in X\), and let \(N = \{x \in X : g(x)h(x) \neq 0\}\), so \(N \in \mathcal{S}\) and \(\mu(N) = 0\). Define

\begin{equation*} A = \{x \in X : g(x) \neq 0\} \quad \text{and} \quad B = \{x \in X : g(x) = 0\}. \end{equation*}

Then \(A, B \in \mathcal{S}\), \(A \cup B = X\), and \(A \cap B = \varnothing\).

If \(x \in A \setminus N\), then \(g(x) \neq 0\) and \(g(x)h(x) = 0\), which forces \(h(x) = 0\). Thus \(h = 0\) on \(A \setminus N\), and since \(\mu(N) = 0\) we get \(\int_{E \cap A} |h| \, d\mu = 0\) for every \(E \in \mathcal{S}\). Hence for every \(E \in \mathcal{S}\),

\begin{equation*} \nu_h(E) - \nu_h(E \cap B) = \int_{E \cap A} h \, d\mu = 0, \end{equation*}

using \(E = (E \cap A) \cup (E \cap B)\) with the two pieces disjoint. So \(\nu_h(E) = \nu_h(E \cap B)\) for all \(E \in \mathcal{S}\).

Also \(g = 0\) everywhere on \(B\) by the definition of \(B\), so for every \(E \in \mathcal{S}\),

\begin{equation*} \nu_g(E) = \int_{E \cap A} g \, d\mu + \int_{E \cap B} g \, d\mu = \int_{E \cap A} g \, d\mu = \nu_g(E \cap A). \end{equation*}

Thus \(A\) and \(B\) satisfy both conditions of 9.28, and therefore \(\nu_g \perp \nu_h\), that is, \(g \, d\mu \perp h \, d\mu\).

Problem (9B.3)

Suppose \(\nu\) and \(\mu\) are complex measures on a measurable space \((X, \mathcal{S})\). Show that the following are equivalent.

(a) \(\nu \perp \mu\).

(b) \(|\nu| \perp |\mu|\).

(c) \(\operatorname{Re} \nu \perp \mu\) and \(\operatorname{Im} \nu \perp \mu\).

Solution

We begin with a lemma that converts the phrase lives on a set, which appears in definition 9.28, into a statement about total variation measures.

Lemma. Suppose \(\sigma\) is a complex measure on \((X, \mathcal{S})\) and \(A, B \in \mathcal{S}\) satisfy \(A \cup B = X\) and \(A \cap B = \varnothing\). Then

\begin{equation*} \sigma(E) = \sigma(E \cap A) \text{ for all } E \in \mathcal{S} \iff |\sigma|(B) = 0. \end{equation*}

Proof of the lemma. Suppose \(\sigma(E) = \sigma(E \cap A)\) for all \(E \in \mathcal{S}\). If \(F \in \mathcal{S}\) and \(F \subseteq B\), then \(F \cap A = \varnothing\) and hence \(\sigma(F) = \sigma(F \cap A) = \sigma(\varnothing) = 0\). Because \(\sigma(F) = 0\) for every measurable \(F \subseteq B\), the third bullet point following 9.8 gives \(|\sigma|(B) = 0\). Conversely, suppose \(|\sigma|(B) = 0\). If \(E \in \mathcal{S}\), then \(|\sigma(E \cap B)| \leq |\sigma|(E \cap B) \leq |\sigma|(B) = 0\) (the first inequality is the first bullet point following 9.8, the second holds because \(|\sigma|\) is a positive measure by 9.11). Hence \(\sigma(E \cap B) = 0\) and so

\begin{equation*} \sigma(E) = \sigma(E \cap A) + \sigma(E \cap B) = \sigma(E \cap A), \end{equation*}

completing the proof of the lemma. Note that the same equivalence holds verbatim when \(\sigma\) is a finite positive measure, with \(|\sigma|\) replaced by \(\sigma\): indeed \(\sigma(E) = \sigma(E \cap A)\) for all \(E\) forces \(\sigma(B) = \sigma(B \cap A) = 0\), and \(\sigma(B) = 0\) forces \(\sigma(E \cap B) = 0\) and hence \(\sigma(E) = \sigma(E \cap A)\).

By the lemma and definition 9.28, for complex measures \(\nu, \mu\) we have

\begin{equation*} \nu \perp \mu \iff \text{there exist disjoint } A, B \in \mathcal{S} \text{ with } A \cup B = X, \ |\nu|(B) = 0, \ |\mu|(A) = 0. \end{equation*}

We refer to the equivalence displayed above as the singularity criterion.

(a) \(\iff\) (b). Because \(|\nu|\) and \(|\mu|\) are finite positive measures (by 9.11 and 9.17), the positive-measure form of the lemma shows that \(|\nu| \perp |\mu|\) holds if and only if there exist disjoint \(A, B \in \mathcal{S}\) with \(A \cup B = X\), \(|\nu|(B) = 0\), and \(|\mu|(A) = 0\). That is exactly the right side of the singularity criterion, which is equivalent to \(\nu \perp \mu\). Hence (a) and (b) are equivalent.

(a) \(\implies\) (c). Suppose \(\nu \perp \mu\), and let \(A, B \in \mathcal{S}\) be as in 9.28, so \(A \cup B = X\), \(A \cap B = \varnothing\), and \(\nu(E) = \nu(E \cap A)\) and \(\mu(E) = \mu(E \cap B)\) for all \(E \in \mathcal{S}\). Taking real parts of the first equation gives

\begin{equation*} (\operatorname{Re}\nu)(E) = \operatorname{Re}\bigl(\nu(E)\bigr) = \operatorname{Re}\bigl(\nu(E \cap A)\bigr) = (\operatorname{Re}\nu)(E \cap A) \end{equation*}

for all \(E \in \mathcal{S}\), and taking imaginary parts gives \((\operatorname{Im}\nu)(E) = (\operatorname{Im}\nu)(E \cap A)\) for all \(E \in \mathcal{S}\). Thus the same pair \(A, B\) witnesses \(\operatorname{Re}\nu \perp \mu\) and \(\operatorname{Im}\nu \perp \mu\), proving (c).

(c) \(\implies\) (a). Suppose \(\operatorname{Re}\nu \perp \mu\) and \(\operatorname{Im}\nu \perp \mu\). By the singularity criterion applied to each of these two relations, there are disjoint pairs \(A_1, B_1\) and \(A_2, B_2\), each with union \(X\), such that

\begin{equation*} |\operatorname{Re}\nu|(B_1) = 0, \quad |\mu|(A_1) = 0, \quad |\operatorname{Im}\nu|(B_2) = 0, \quad |\mu|(A_2) = 0. \end{equation*}

Set

\begin{equation*} A = A_1 \cup A_2 \quad \text{and} \quad B = X \setminus A = B_1 \cap B_2 . \end{equation*}

Then \(A, B \in \mathcal{S}\) are disjoint with union \(X\). Because \(|\mu|\) is a positive measure,

\begin{equation*} |\mu|(A) \leq |\mu|(A_1) + |\mu|(A_2) = 0 . \end{equation*}

Next we bound \(|\nu|(B)\). If \(n \in \mathbf{Z}^+\) and \(E_1, \dots, E_n\) are disjoint sets in \(\mathcal{S}\) with \(E_1 \cup \cdots \cup E_n \subseteq B\), then

\begin{equation*} \sum_{k=1}^n |\nu(E_k)| \leq \sum_{k=1}^n |(\operatorname{Re}\nu)(E_k)| + \sum_{k=1}^n |(\operatorname{Im}\nu)(E_k)| \leq |\operatorname{Re}\nu|(B) + |\operatorname{Im}\nu|(B), \end{equation*}

where the first inequality uses \(|z| \leq |\operatorname{Re} z| + |\operatorname{Im} z|\) and the second uses definition 9.8. Taking the supremum over all such collections gives

\begin{equation*} |\nu|(B) \leq |\operatorname{Re}\nu|(B) + |\operatorname{Im}\nu|(B) \leq |\operatorname{Re}\nu|(B_1) + |\operatorname{Im}\nu|(B_2) = 0, \end{equation*}

where the second inequality holds because \(B \subseteq B_1\), \(B \subseteq B_2\), and the total variation measures are positive measures. Hence \(|\nu|(B) = 0\) and \(|\mu|(A) = 0\), so the singularity criterion gives \(\nu \perp \mu\), proving (a).

Thus (a), (b), and (c) are equivalent.

Problem (9B.4)

Suppose \(\nu\) and \(\mu\) are complex measures on a measurable space \((X, \mathcal{S})\). Prove that if \(\nu \perp \mu\), then \(|\nu + \mu| = |\nu| + |\mu|\) and \(\|\nu + \mu\| = \|\nu\| + \|\mu\|\).

Solution

Suppose \(\nu \perp \mu\). By 9.28 there exist \(A, B \in \mathcal{S}\) such that

\begin{equation*} A \cup B = X, \quad A \cap B = \varnothing, \quad \nu(E) = \nu(E \cap A), \quad \mu(E) = \mu(E \cap B) \end{equation*}

for all \(E \in \mathcal{S}\). These conditions imply

\begin{equation*} |\nu|(B) = 0 \quad \text{and} \quad |\mu|(A) = 0 . \end{equation*}

Indeed, if \(F \in \mathcal{S}\) and \(F \subseteq B\), then \(F \cap A = \varnothing\) and hence \(\nu(F) = \nu(F \cap A) = \nu(\varnothing) = 0\); because \(\nu(F) = 0\) for every measurable \(F \subseteq B\), the third bullet point following 9.8 gives \(|\nu|(B) = 0\). Symmetrically, every measurable \(F \subseteq A\) satisfies \(\mu(F) = \mu(F \cap B) = 0\), so \(|\mu|(A) = 0\).

We also record an evident consequence of definition 9.8: if \(\sigma\) and \(\tau\) are complex measures on \((X, \mathcal{S})\), \(C \in \mathcal{S}\), and \(\sigma(F) = \tau(F)\) for every \(F \in \mathcal{S}\) with \(F \subseteq C\), then

\begin{equation*} |\sigma|(D) = |\tau|(D) \quad \text{for every } D \in \mathcal{S} \text{ with } D \subseteq C, \end{equation*}

because the supremum defining \(|\sigma|(D)\) ranges over finite disjoint collections \(E_1, \dots, E_n\) in \(\mathcal{S}\) with \(E_1 \cup \cdots \cup E_n \subseteq D \subseteq C\), and for such a collection each \(E_k \subseteq C\), so \(\sigma(E_k) = \tau(E_k)\) and the two suprema are suprema of the same set of numbers.

Now fix \(E \in \mathcal{S}\). Because \(|\mu|(A) = 0\), every \(F \in \mathcal{S}\) with \(F \subseteq A\) satisfies \(|\mu(F)| \leq |\mu|(F) \leq |\mu|(A) = 0\), hence \(\mu(F) = 0\) and therefore

\begin{equation*} (\nu + \mu)(F) = \nu(F) \quad \text{for every } F \in \mathcal{S} \text{ with } F \subseteq A . \end{equation*}

By the observation of the previous paragraph (with \(C = A\)),

\begin{equation*} |\nu + \mu|(E \cap A) = |\nu|(E \cap A). \end{equation*}

Symmetrically, \(|\nu|(B) = 0\) gives \((\nu + \mu)(F) = \mu(F)\) for every measurable \(F \subseteq B\), and hence

\begin{equation*} |\nu + \mu|(E \cap B) = |\mu|(E \cap B). \end{equation*}

Furthermore \(|\nu|(E \cap B) \leq |\nu|(B) = 0\) and \(|\mu|(E \cap A) \leq |\mu|(A) = 0\), so by the additivity of the positive measures \(|\nu|\) and \(|\mu|\) (9.11),

\begin{equation*} |\nu|(E) = |\nu|(E \cap A) + |\nu|(E \cap B) = |\nu|(E \cap A) \end{equation*}

and

\begin{equation*} |\mu|(E) = |\mu|(E \cap A) + |\mu|(E \cap B) = |\mu|(E \cap B). \end{equation*}

Putting the pieces together and using the additivity of the positive measure \(|\nu + \mu|\),

\begin{equation*} \begin{aligned} |\nu + \mu|(E) &= |\nu + \mu|(E \cap A) + |\nu + \mu|(E \cap B) \\ &= |\nu|(E \cap A) + |\mu|(E \cap B) \\ &= |\nu|(E) + |\mu|(E). \end{aligned} \end{equation*}

Because \(E \in \mathcal{S}\) was arbitrary, this proves that \(|\nu + \mu| = |\nu| + |\mu|\) as measures on \((X, \mathcal{S})\).

Finally, taking \(E = X\) in the equation above and using definition 9.15 of the total variation norm gives

\begin{equation*} \|\nu + \mu\| = |\nu + \mu|(X) = |\nu|(X) + |\mu|(X) = \|\nu\| + \|\mu\| , \end{equation*}

as desired.

Problem (9B.5)

Suppose \(\nu\) and \(\mu\) are finite (positive) measures on a measurable space \((X, \mathcal{S})\). Prove that \(\nu \perp \mu\) if and only if \(\|\nu - \mu\| = \|\nu\| + \|\mu\|\).

Solution

Throughout, recall from the first bullet point of 9.16 that a finite positive measure \(\sigma\) satisfies \(\|\sigma\| = \sigma(X)\). In particular \(\|\nu\| = \nu(X)\) and \(\|\mu\| = \mu(X)\), and both are finite.

First suppose \(\nu \perp \mu\). Let \(A, B \in \mathcal{S}\) be as in 9.28, so \(A \cup B = X\), \(A \cap B = \varnothing\), and \(\nu(E) = \nu(E \cap A)\) and \(\mu(E) = \mu(E \cap B)\) for all \(E \in \mathcal{S}\). Taking \(E = B\) in the first of these equations gives \(\nu(B) = \nu(B \cap A) = \nu(\varnothing) = 0\), and taking \(E = A\) in the second gives \(\mu(A) = \mu(A \cap B) = 0\). Hence, by additivity,

\begin{equation*} \nu(A) = \nu(X) - \nu(B) = \nu(X) \quad \text{and} \quad \mu(B) = \mu(X) - \mu(A) = \mu(X), \end{equation*}

where the subtractions are legitimate because \(\nu\) and \(\mu\) are finite.

Let \(\sigma = \nu - \mu\), a real measure on \((X, \mathcal{S})\). Because \(A, B\) are disjoint sets in \(\mathcal{S}\) with \(A \cup B = X\), definition 9.8 gives

\begin{equation*} \|\sigma\| = |\sigma|(X) \geq |\sigma(A)| + |\sigma(B)| = |\nu(A) - \mu(A)| + |\nu(B) - \mu(B)| = \nu(X) + \mu(X) = \|\nu\| + \|\mu\| . \end{equation*}

For the reverse inequality, suppose \(n \in \mathbf{Z}^+\) and \(E_1, \dots, E_n\) are disjoint sets in \(\mathcal{S}\). Then

\begin{equation*} \sum_{k=1}^n |\sigma(E_k)| \leq \sum_{k=1}^n |\nu(E_k)| + \sum_{k=1}^n |\mu(E_k)| \leq |\nu|(X) + |\mu|(X) = \|\nu\| + \|\mu\| , \end{equation*}

where the second inequality uses definition 9.8 twice. Taking the supremum over all such collections gives \(\|\sigma\| \leq \|\nu\| + \|\mu\|\). Combining the two inequalities,

\begin{equation*} \|\nu - \mu\| = \|\nu\| + \|\mu\| . \end{equation*}

Conversely, suppose \(\|\nu - \mu\| = \|\nu\| + \|\mu\|\). Let \(\sigma = \nu - \mu\), which is a real measure on \((X, \mathcal{S})\) because \(\nu\) and \(\mu\) are finite positive measures. By the Hahn Decomposition Theorem (9.23) there exist \(A, B \in \mathcal{S}\) with \(A \cup B = X\), \(A \cap B = \varnothing\), \(\sigma(E) \geq 0\) for all measurable \(E \subseteq A\), and \(\sigma(E) \leq 0\) for all measurable \(E \subseteq B\).

As in the proof of the Jordan Decomposition Theorem (9.30), define finite positive measures \(\sigma^+\) and \(\sigma^-\) by \(\sigma^+(E) = \sigma(E \cap A)\) and \(\sigma^-(E) = -\sigma(E \cap B)\); then \(|\sigma| = \sigma^+ + \sigma^-\) by 9.31. Hence

\begin{equation*} \|\sigma\| = |\sigma|(X) = \sigma^+(X) + \sigma^-(X) = \sigma(A) - \sigma(B). \end{equation*}

Writing this out in terms of \(\nu\) and \(\mu\),

\begin{equation*} \|\nu - \mu\| = \bigl(\nu(A) - \mu(A)\bigr) - \bigl(\nu(B) - \mu(B)\bigr) = \nu(A) - \nu(B) - \mu(A) + \mu(B). \end{equation*}

On the other hand, since \(A\) and \(B\) are disjoint with union \(X\),

\begin{equation*} \|\nu\| + \|\mu\| = \nu(X) + \mu(X) = \nu(A) + \nu(B) + \mu(A) + \mu(B). \end{equation*}

All the numbers appearing here are finite, so equating the two right sides and cancelling \(\nu(A)\) and \(\mu(B)\) gives

\begin{equation*}

  • \nu(B) - \mu(A) = \nu(B) + \mu(A), \end{equation*}

which yields

\begin{equation*} 2\bigl(\nu(B) + \mu(A)\bigr) = 0 . \end{equation*}

Because \(\nu\) and \(\mu\) are positive measures, \(\nu(B) \geq 0\) and \(\mu(A) \geq 0\), so

\begin{equation*} \nu(B) = 0 \quad \text{and} \quad \mu(A) = 0 . \end{equation*}

Now for every \(E \in \mathcal{S}\) we have \(\nu(E \cap B) \leq \nu(B) = 0\) and \(\mu(E \cap A) \leq \mu(A) = 0\), so

\begin{equation*} \nu(E) = \nu(E \cap A) + \nu(E \cap B) = \nu(E \cap A) \end{equation*}

and

\begin{equation*} \mu(E) = \mu(E \cap A) + \mu(E \cap B) = \mu(E \cap B). \end{equation*}

Thus \(A\) and \(B\) satisfy the conditions of 9.28, and therefore \(\nu \perp \mu\).

Problem (9B.6)

Suppose \(\mu\) is a complex or positive measure on a measurable space \((X, \mathcal{S})\). Prove that

\begin{equation*} \{\nu \in \mathcal{M}_{\mathbf{F}}(\mathcal{S}) : \nu \perp \mu\} \end{equation*}

is a closed subspace of \(\mathcal{M}_{\mathbf{F}}(\mathcal{S})\).

Solution

Let \(V = \{\nu \in \mathcal{M}_{\mathbf{F}}(\mathcal{S}) : \nu \perp \mu\}\). Call a pair \(A, B \in \mathcal{S}\) a partition of \(X\) if \(A \cup B = X\) and \(A \cap B = \varnothing\).

We first restate the two conditions in definition 9.28 in a form convenient for the proof. Note that \(\mu\) is allowed to be a positive measure that is not finite, so we avoid using \(|\mu|\).

Observation 1. Suppose \(\sigma\) is a complex or positive measure on \((X, \mathcal{S})\) and \(A, B\) is a partition of \(X\). Then

\begin{equation*} \sigma(E) = \sigma(E \cap B) \text{ for all } E \in \mathcal{S} \iff \sigma(F) = 0 \text{ for all } F \in \mathcal{S} \text{ with } F \subseteq A . \end{equation*}

Indeed, if the left condition holds and \(F \in \mathcal{S}\) with \(F \subseteq A\), then \(\sigma(F) = \sigma(F \cap B) = \sigma(\varnothing) = 0\). Conversely, if the right condition holds and \(E \in \mathcal{S}\), then \(E \cap A \subseteq A\), so \(\sigma(E \cap A) = 0\) and additivity gives \(\sigma(E) = \sigma(E \cap A) + \sigma(E \cap B) = \sigma(E \cap B)\).

Observation 2. Suppose \(\nu \in \mathcal{M}_{\mathbf{F}}(\mathcal{S})\) and \(A, B\) is a partition of \(X\). Then, by Observation 1 with the roles of \(A\) and \(B\) interchanged together with the third bullet point following 9.8,

\begin{equation*} \nu(E) = \nu(E \cap A) \text{ for all } E \in \mathcal{S} \iff \nu(F) = 0 \text{ for all measurable } F \subseteq B \iff |\nu|(B) = 0 . \end{equation*}

Combining the two observations with 9.28: for \(\nu \in \mathcal{M}_{\mathbf{F}}(\mathcal{S})\),

\begin{equation*} \nu \perp \mu \iff \text{there is a partition } A, B \text{ of } X \text{ with } |\nu|(B) = 0 \text{ and } \mu(F) = 0 \text{ for all measurable } F \subseteq A. \end{equation*}

We refer to the equivalence displayed above as the singularity criterion.

A partition \(A, B\) of \(X\) with the two properties on the right side of that criterion will be called a witness for \(\nu\).

\(V\) contains \(0\). Take \(A = \varnothing\) and \(B = X\). Then \(|0|(B) = 0\), and the condition on \(\mu\) is vacuous since the only measurable subset of \(\varnothing\) is \(\varnothing\), which has \(\mu(\varnothing) = 0\). So \(0 \in V\).

\(V\) is closed under scalar multiplication. Suppose \(\nu \in V\) with witness \(A, B\), and \(\alpha \in \mathbf{F}\). Directly from definition 9.8, \(|\alpha \nu| = |\alpha| \, |\nu|\), so \(|\alpha\nu|(B) = |\alpha| \cdot 0 = 0\). Thus \(A, B\) is also a witness for \(\alpha\nu\), so \(\alpha \nu \in V\).

\(V\) is closed under addition. Suppose \(\nu_1, \nu_2 \in V\), with witnesses \(A_1, B_1\) and \(A_2, B_2\). Let

\begin{equation*} A = A_1 \cup A_2 \quad \text{and} \quad B = X \setminus A = B_1 \cap B_2, \end{equation*}

so \(A, B\) is a partition of \(X\). If \(F \in \mathcal{S}\) and \(F \subseteq A\), write \(F = (F \cap A_1) \cup (F \setminus A_1)\); these are disjoint measurable sets with \(F \cap A_1 \subseteq A_1\) and \(F \setminus A_1 \subseteq A_2\), so \(\mu(F) = \mu(F \cap A_1) + \mu(F \setminus A_1) = 0 + 0 = 0\). Next, \(B \subseteq B_1\) and \(B \subseteq B_2\), so monotonicity of the positive measures \(|\nu_1|\) and \(|\nu_2|\) (9.11) gives \(|\nu_1|(B) = |\nu_2|(B) = 0\). Finally, directly from definition 9.8, if \(E_1, \dots, E_n\) are disjoint sets in \(\mathcal{S}\) with union contained in \(B\), then

\begin{equation*} \sum_{k=1}^n |(\nu_1 + \nu_2)(E_k)| \leq \sum_{k=1}^n |\nu_1(E_k)| + \sum_{k=1}^n |\nu_2(E_k)| \leq |\nu_1|(B) + |\nu_2|(B) = 0, \end{equation*}

so \(|\nu_1 + \nu_2|(B) = 0\). Hence \(A, B\) is a witness for \(\nu_1 + \nu_2\), so \(\nu_1 + \nu_2 \in V\). Therefore \(V\) is a subspace of \(\mathcal{M}_{\mathbf{F}}(\mathcal{S})\).

\(V\) is closed. Suppose \(\nu_1, \nu_2, \dots\) is a sequence in \(V\) and \(\nu \in \mathcal{M}_{\mathbf{F}}(\mathcal{S})\) with \(\lim_{k \to \infty} \|\nu - \nu_k\| = 0\). For each \(k \in \mathbf{Z}^+\) let \(A_k, B_k\) be a witness for \(\nu_k\). Let

\begin{equation*} A = \bigcup_{k=1}^\infty A_k \quad \text{and} \quad B = X \setminus A = \bigcap_{k=1}^\infty B_k , \end{equation*}

which is a partition of \(X\) with \(A, B \in \mathcal{S}\).

First we check the condition on \(\mu\). Suppose \(F \in \mathcal{S}\) and \(F \subseteq A\). Let \(F_1 = F \cap A_1\) and, for \(k \geq 2\), let

\begin{equation*} F_k = (F \cap A_k) \setminus (A_1 \cup \cdots \cup A_{k-1}). \end{equation*}

Then \(F_1, F_2, \dots\) are disjoint sets in \(\mathcal{S}\) with \(F_k \subseteq A_k\) and \(\bigcup_{k=1}^\infty F_k = F\). Because \(A_k, B_k\) is a witness for \(\nu_k\), we have \(\mu(F_k) = 0\) for each \(k\). Countable additivity of \(\mu\) (valid whether \(\mu\) is complex or positive) gives

\begin{equation*} \mu(F) = \sum_{k=1}^\infty \mu(F_k) = 0 . \end{equation*}

Next we check that \(|\nu|(B) = 0\). Fix \(k \in \mathbf{Z}^+\). Because \(\nu = (\nu - \nu_k) + \nu_k\), the same argument as in the additivity step above (using definition 9.8) gives

\begin{equation*} |\nu|(B) \leq |\nu - \nu_k|(B) + |\nu_k|(B). \end{equation*}

Now \(B \subseteq B_k\), so \(|\nu_k|(B) \leq |\nu_k|(B_k) = 0\), and \(|\nu - \nu_k|(B) \leq |\nu - \nu_k|(X) = \|\nu - \nu_k\|\). Hence

\begin{equation*} 0 \leq |\nu|(B) \leq \|\nu - \nu_k\| \quad \text{for every } k \in \mathbf{Z}^+ . \end{equation*}

Letting \(k \to \infty\) gives \(|\nu|(B) = 0\).

Thus \(A, B\) is a witness for \(\nu\), so the singularity criterion gives \(\nu \perp \mu\), that is, \(\nu \in V\). Hence \(V\) contains the limit of every convergent sequence of its elements, so \(V\) is closed.

Therefore \(V = \{\nu \in \mathcal{M}_{\mathbf{F}}(\mathcal{S}) : \nu \perp \mu\}\) is a closed subspace of \(\mathcal{M}_{\mathbf{F}}(\mathcal{S})\). Because \(\mathcal{M}_{\mathbf{F}}(\mathcal{S})\) is a Banach space (9.18), \(V\) is itself a Banach space with the total variation norm.

Problem (9B.7)

Use the Cantor set to prove that there exists a (positive) measure \(\nu\) on \((\mathbf{R}, \mathcal{B})\) such that \(\nu \perp \lambda\) and \(\nu(\mathbf{R}) \neq 0\) but \(\nu(\{x\}) = 0\) for every \(x \in \mathbf{R}\); here \(\lambda\) denotes Lebesgue measure on the \(\sigma\)-algebra \(\mathcal{B}\) of Borel subsets of \(\mathbf{R}\).

[The second bullet point in Example 9.29 does not provide an example of the desired behavior because in that example, \(\nu(\{r_k\}) \neq 0\) for all \(k \in \mathbf{Z}^+\) with \(w_k \neq 0\).]

Solution

Let \(C\) denote the Cantor set (2.74) and let \(\Lambda : [0,1] \to [0,1]\) denote the Cantor function (2.77). By 2.79, \(\Lambda\) is continuous, increasing, and maps \([0,1]\) onto \([0,1]\). By 2.76(b), \(\lambda( C) = 0\); note also that \(C\) is closed (2.76(a)) and hence a Borel set.

We build the desired measure as the pushforward of Lebesgue measure on \([0,1]\) under a right inverse of \(\Lambda\) that takes values in \(C\).

Definition of the map. For \(y \in [0,1]\) let

\begin{equation*} S_y = \{x \in [0,1] : \Lambda(x) \geq y\} \quad \text{and} \quad f(y) = \inf S_y . \end{equation*}

The set \(S_y\) is nonempty because \(\Lambda(1) = 1 \geq y\), so \(f(y) \in [0,1]\) is well defined. This defines \(f : [0,1] \to [0,1]\).

Step 1: \(\Lambda(f(y)) = y\) for every \(y \in [0,1]\). Because \(\Lambda\) is continuous, \(S_y = \Lambda^{-1}\bigl([y, \infty)\bigr)\) is a closed subset of \([0,1]\), so \(f(y) = \inf S_y \in S_y\); that is, \(\Lambda(f(y)) \geq y\). For the reverse inequality, first suppose \(f(y) = 0\). Then \(0 = \Lambda(0) = \Lambda(f(y)) \geq y \geq 0\), so \(\Lambda(f(y)) = 0 = y\). Now suppose \(f(y) > 0\). If \(0 \leq x < f(y)\), then \(x \notin S_y\) and hence \(\Lambda(x) < y\). Letting \(x\) increase to \(f(y)\) and using the continuity of \(\Lambda\) gives \(\Lambda(f(y)) \leq y\). In both cases \(\Lambda(f(y)) = y\).

Step 2: \(f\) is increasing and injective. If \(0 \leq y \leq y^{\prime} \leq 1\), then \(S_{y^{\prime}} \subseteq S_y\), so \(f(y) = \inf S_y \leq \inf S_{y^{\prime}} = f(y^{\prime})\); thus \(f\) is increasing. If \(f(y) = f(y^{\prime})\), then applying \(\Lambda\) and using Step 1 gives \(y = \Lambda(f(y)) = \Lambda(f(y^{\prime})) = y^{\prime}\); thus \(f\) is injective.

Step 3: \(f(y) \in C\) for every \(y \in [0,1]\). First note that \(0 \in C\) and \(1 \in C\): by 2.74 each open interval removed in the construction of the Cantor set is the middle third \(\bigl(\tfrac{2c+d}{3}, \tfrac{c+2d}{3}\bigr)\) of some closed subinterval \([c,d] \subseteq [0,1]\), and this middle third is contained in \((c,d) \subseteq (0,1)\). Hence no removed interval contains \(0\) or \(1\), so \(0\) and \(1\) survive in \(C\).

Now suppose, seeking a contradiction, that \(f(y) \in [0,1] \setminus C\). By the previous paragraph, \(0 < f(y) < 1\). As noted in the proof of 2.79, \(\Lambda\) is constant on some open interval containing \(f(y)\); intersecting that interval with \((0,1)\), we obtain \(a, b\) with \(0 \leq a < f(y) < b \leq 1\) such that \(\Lambda\) is constant on \((a,b)\), say with value \(c\). Because \(f(y) \in S_y\) (established in Step 1), we have \(c = \Lambda(f(y)) \geq y\). On the other hand, choose \(x\) with \(a < x < f(y)\); then \(x \in (0,1) \subseteq [0,1]\) and \(x < f(y) = \inf S_y\), so \(x \notin S_y\) and hence \(\Lambda(x) < y\). But \(x \in (a,b)\), so \(\Lambda(x) = c\). Thus \(c < y\), contradicting \(c \geq y\). Hence \(f(y) \in C\).

Step 4: \(f\) is Borel measurable. The domain \([0,1]\) is a Borel subset of \(\mathbf{R}\) and \(f\) is increasing by Step 2, so \(f\) is a Borel measurable function by 2.43. By the definition 2.40 of Borel measurability, this means precisely that \(f^{-1}(E)\) is a Borel subset of \(\mathbf{R}\) for every \(E \in \mathcal{B}\); note also that \(f^{-1}(E) \subseteq [0,1]\).

Definition of \(\nu\). Define \(\nu : \mathcal{B} \to [0, \infty)\) by

\begin{equation*} \nu(E) = \lambda\bigl(f^{-1}(E)\bigr) \quad \text{for } E \in \mathcal{B} , \end{equation*}

which makes sense by Step 4. Then \(\nu(\varnothing) = \lambda(\varnothing) = 0\). If \(E_1, E_2, \dots\) is a disjoint sequence in \(\mathcal{B}\), then \(f^{-1}(E_1), f^{-1}(E_2), \dots\) are disjoint Borel sets whose union is \(f^{-1}\bigl(\bigcup_{k=1}^\infty E_k\bigr)\), so the countable additivity of \(\lambda\) gives

\begin{equation*} \nu\Bigl( \bigcup_{k=1}^\infty E_k \Bigr) = \lambda\Bigl( \bigcup_{k=1}^\infty f^{-1}(E_k) \Bigr) = \sum_{k=1}^\infty \lambda\bigl(f^{-1}(E_k)\bigr) = \sum_{k=1}^\infty \nu(E_k). \end{equation*}

Thus \(\nu\) is a (positive) measure on \((\mathbf{R}, \mathcal{B})\).

\(\nu(\mathbf{R}) \neq 0\). Because \(f\) is defined on all of \([0,1]\), we have \(f^{-1}(\mathbf{R}) = [0,1]\), so

\begin{equation*} \nu(\mathbf{R}) = \lambda([0,1]) = 1 \neq 0 . \end{equation*}

\(\nu \perp \lambda\). By Step 3 the range of \(f\) is contained in \(C\), so \(f^{-1}(\mathbf{R} \setminus C) = \varnothing\) and hence \(\nu(\mathbf{R} \setminus C) = 0\). Let \(A = C\) and \(B = \mathbf{R} \setminus C\); these are disjoint Borel sets with union \(\mathbf{R}\). For every \(E \in \mathcal{B}\),

\begin{equation*} \nu(E) = \nu(E \cap A) + \nu(E \cap B) = \nu(E \cap A), \end{equation*}

because \(\nu(E \cap B) \leq \nu(\mathbf{R} \setminus C) = 0\). Also, since \(\lambda( C) = 0\) by 2.76(b),

\begin{equation*} \lambda(E) = \lambda(E \cap A) + \lambda(E \cap B) = \lambda(E \cap B), \end{equation*}

because \(\lambda(E \cap A) \leq \lambda( C) = 0\). Thus \(A\) and \(B\) satisfy the conditions of 9.28, and therefore \(\nu \perp \lambda\).

\(\nu(\{x\}) = 0\) for every \(x \in \mathbf{R}\). Fix \(x \in \mathbf{R}\). Because \(f\) is injective (Step 2), the set \(f^{-1}(\{x\})\) contains at most one point of \([0,1]\). A set containing at most one point has Lebesgue measure \(0\), so

\begin{equation*} \nu(\{x\}) = \lambda\bigl(f^{-1}(\{x\})\bigr) = 0 . \end{equation*}

Hence \(\nu\) is a positive measure on \((\mathbf{R}, \mathcal{B})\) with \(\nu \perp \lambda\), \(\nu(\mathbf{R}) = 1 \neq 0\), and \(\nu(\{x\}) = 0\) for every \(x \in \mathbf{R}\), as desired.

Remark. The measure \(\nu\) constructed above is the classical Cantor measure. It lives entirely on the Cantor set, which has Lebesgue measure \(0\), yet it gives no mass to any individual point; this is possible precisely because the Cantor set is uncountable (2.80). Equivalently, \(\nu\) is the measure whose distribution function is the Cantor function. Indeed, Step 1 and Step 2 give \(\{y \in [0,1] : f(y) \leq t\} = [0, \Lambda(t)]\) for each \(t \in [0,1]\) (if \(f(y) \leq t\) then \(y = \Lambda(f(y)) \leq \Lambda(t)\) because \(\Lambda\) is increasing, and if \(y \leq \Lambda(t)\) then \(t \in S_y\) and so \(f(y) \leq t\)), and since \(\{y : f(y) < s\}\) differs from \([0, \Lambda(s)]\) by at most one point, we get \(\nu([s,t]) = \Lambda(t) - \Lambda(s)\) for \(0 \leq s \leq t \leq 1\). The continuity of \(\Lambda\) is exactly what makes every singleton \(\nu\)-null.

Problem (9B.8)

Suppose \(\nu\) is a real measure on a measurable space \((X, \mathcal{S})\). Prove that

\begin{equation*} \nu^{+}(E) = \sup\{\nu(D) : D \in \mathcal{S} \text{ and } D \subseteq E\} \end{equation*}

and

\begin{equation*} \nu^{-}(E) = -\inf\{\nu(D) : D \in \mathcal{S} \text{ and } D \subseteq E\} \end{equation*}

for all \(E \in \mathcal{S}\).

Solution

Let \(A, B \in \mathcal{S}\) be a Hahn decomposition of \(\nu\) as in 9.23: thus \(A \cup B = X\), \(A \cap B = \varnothing\), \(\nu(F) \ge 0\) for every \(F \in \mathcal{S}\) with \(F \subseteq A\), and \(\nu(F) \le 0\) for every \(F \in \mathcal{S}\) with \(F \subseteq B\). The proof of the Jordan Decomposition Theorem (9.30) shows that the measures \(\nu^{+}\) and \(\nu^{-}\) produced by that theorem are given by

\begin{equation*} \nu^{+}(E) = \nu(E \cap A) \qquad \text{and} \qquad \nu^{-}(E) = -\nu(E \cap B) \end{equation*}

for all \(E \in \mathcal{S}\); the uniqueness assertion in 9.30 guarantees that these are indeed the measures denoted \(\nu^{+}\) and \(\nu^{-}\), no matter which Hahn decomposition is used.

Fix \(E \in \mathcal{S}\), and write

\begin{equation*} s = \sup\{\nu(D) : D \in \mathcal{S} \text{ and } D \subseteq E\}, \qquad i = \inf\{\nu(D) : D \in \mathcal{S} \text{ and } D \subseteq E\}. \end{equation*}

Both sets above are nonempty (they contain \(\nu(\varnothing) = 0\), by 9.3(a)), so \(s \ge 0 \ge i\).

First we prove \(\nu^{+}(E) = s\).

The set \(D_{0} = E \cap A\) belongs to \(\mathcal{S}\) and satisfies \(D_{0} \subseteq E\). Hence

\begin{equation*} s \ge \nu(D_{0}) = \nu(E \cap A) = \nu^{+}(E). \end{equation*}

For the reverse inequality, suppose \(D \in \mathcal{S}\) and \(D \subseteq E\). Because \(\nu = \nu^{+} - \nu^{-}\) (by 9.31) and \(\nu^{-}\) is a positive measure,

\begin{equation*} \nu(D) = \nu^{+}(D) - \nu^{-}(D) \le \nu^{+}(D) \le \nu^{+}(E), \end{equation*}

where the last inequality holds because \(\nu^{+}\) is a (positive) measure and \(D \subseteq E\), so that \(\nu^{+}(E) = \nu^{+}(D) + \nu^{+}(E \setminus D) \ge \nu^{+}(D)\). Taking the supremum over all such \(D\) gives \(s \le \nu^{+}(E)\).

Combining the two inequalities yields \(\nu^{+}(E) = s\), which is the first displayed formula. Note that the supremum is attained, at \(D = E \cap A\).

Now we prove \(\nu^{-}(E) = -i\).

The set \(D_{1} = E \cap B\) belongs to \(\mathcal{S}\) and satisfies \(D_{1} \subseteq E\). Hence

\begin{equation*} i \le \nu(D_{1}) = \nu(E \cap B) = -\nu^{-}(E). \end{equation*}

For the reverse inequality, suppose \(D \in \mathcal{S}\) and \(D \subseteq E\). Because \(\nu^{+}\) is a positive measure and \(\nu^{-}\) is a (positive) measure that is monotone,

\begin{equation*} \nu(D) = \nu^{+}(D) - \nu^{-}(D) \ge -\nu^{-}(D) \ge -\nu^{-}(E). \end{equation*}

Taking the infimum over all such \(D\) gives \(i \ge -\nu^{-}(E)\).

Combining the two inequalities yields \(i = -\nu^{-}(E)\), that is, \(\nu^{-}(E) = -i\), which is the second displayed formula. Again the infimum is attained, at \(D = E \cap B\).

Problem (9B.9)

Suppose \(\mu\) is a (positive) finite measure on a measurable space \((X, \mathcal{S})\) and \(h\) is a nonnegative function in \(L^{1}(\mu)\). Thus \(h \, d\mu \ll d\mu\). Find a reasonable condition on \(h\) that is equivalent to the condition \(d\mu \ll h \, d\mu\).

Solution

The condition is

\begin{equation*} \mu(\{x \in X : h(x) = 0\}) = 0, \end{equation*}

in other words, \(h(x) > 0\) for almost every \(x \in X\).

Throughout, let \(\nu\) denote the finite (positive) measure defined by \(d\nu = h \, d\mu\), so

\begin{equation*} \nu(E) = \int_{E} h \, d\mu \qquad \text{for } E \in \mathcal{S}, \end{equation*}

and let

\begin{equation*} Z = \{x \in X : h(x) = 0\}, \end{equation*}

which belongs to \(\mathcal{S}\) because \(h\) is \(\mathcal{S}\)-measurable. Note that \(\mu\) is a finite (positive) measure and hence is a real measure, so the statement \(\mu \ll \nu\) makes sense in the sense of 9.32.

First suppose \(\mu(Z) = 0\); we show \(\mu \ll \nu\). Suppose \(E \in \mathcal{S}\) and \(\nu(E) = 0\), that is, \(\int_{E} h \, d\mu = 0\). For each \(n \in \mathbf{Z}^{+}\) let

\begin{equation*} E_{n} = \Bigl\{x \in E : h(x) > \tfrac{1}{n}\Bigr\}. \end{equation*}

Then \(\tfrac{1}{n}\chi_{E_{n}} \le h\chi_{E}\) pointwise on \(X\), and hence

\begin{equation*} \tfrac{1}{n}\,\mu(E_{n}) = \int \tfrac{1}{n}\chi_{E_{n}} \, d\mu \le \int_{E} h \, d\mu = 0, \end{equation*}

so \(\mu(E_{n}) = 0\). Because \(h \ge 0\), we have \(\{x \in E : h(x) > 0\} = E \setminus Z\), and hence \(E_{1} \subseteq E_{2} \subseteq \cdots\) with \(\bigcup_{n=1}^{\infty} E_{n} = E \setminus Z\). Thus 2.59 (the measure of an increasing union) gives

\begin{equation*} \mu(E \setminus Z) = \lim_{n \to \infty} \mu(E_{n}) = 0. \end{equation*}

Therefore

\begin{equation*} \mu(E) = \mu(E \setminus Z) + \mu(E \cap Z) \le 0 + \mu(Z) = 0. \end{equation*}

Thus \(\mu(E) = 0\), proving that \(\mu \ll \nu\).

Conversely suppose \(\mu \ll \nu\). Because \(h = 0\) on \(Z\), we have

\begin{equation*} \nu(Z) = \int_{Z} h \, d\mu = 0, \end{equation*}

and hence \(\mu(Z) = 0\) by the definition of absolute continuity. This is exactly the stated condition.

Thus \(d\mu \ll h \, d\mu\) if and only if \(h > 0\) almost everywhere with respect to \(\mu\).

Two remarks. First, the hypothesis that \(\mu\) is finite was used only to make \(\mu\) a real measure, so that the relation \(\mu \ll \nu\) falls under 9.32; the equivalence itself is valid for an arbitrary (positive) measure \(\mu\) if one reads \(\mu \ll \nu\) as the statement that \(\nu(E) = 0\) implies \(\mu(E) = 0\).

Second, the condition \(h > 0\) almost everywhere is exactly what one expects, because it is the condition under which \(1/h\) can serve as a density for \(\mu\) with respect to \(\nu\). Indeed, define \(g : X \to [0, \infty)\) by \(g(x) = 1/h(x)\) if \(h(x) > 0\) and \(g(x) = 0\) if \(h(x) = 0\). The identity \(\int f \, d\nu = \int f h \, d\mu\) holds for \(f = \chi_{E}\) by the definition of \(\nu\), hence for nonnegative simple \(f\) by linearity, hence for all nonnegative \(\mathcal{S}\)-measurable \(f\) by the Monotone Convergence Theorem (3.11). Applying it with \(f = g\chi_{E}\) gives

\begin{equation*} \int_{E} g \, d\nu = \int_{E} gh \, d\mu = \int_{E} \chi_{X \setminus Z} \, d\mu = \mu(E \setminus Z) = \mu(E) \end{equation*}

for every \(E \in \mathcal{S}\), where the last equality uses \(\mu(Z) = 0\). In particular \(g \in L^{1}(\nu)\) because \(\int_{X} g \, d\nu = \mu(X) < \infty\), so \(d\mu = g \, d\nu\), in agreement with the Radon–Nikodym Theorem (9.36).

Problem (9B.10)

Suppose \(\mu\) is a (positive) measure on a measurable space \((X, \mathcal{S})\) and \(\nu\) is a complex measure on \((X, \mathcal{S})\). Show that the following are equivalent.

(a)
\(\nu \ll \mu\).
(b)
\(|\nu| \ll \mu\).
(c)
\(\operatorname{Re} \nu \ll \mu\) and \(\operatorname{Im} \nu \ll \mu\).
Solution

Recall that \(\operatorname{Re}\nu\) and \(\operatorname{Im}\nu\) are the real measures on \((X, \mathcal{S})\) defined by \((\operatorname{Re}\nu)(E) = \operatorname{Re}\bigl(\nu(E)\bigr)\) and \((\operatorname{Im}\nu)(E) = \operatorname{Im}\bigl(\nu(E)\bigr)\), and that \(|\nu|\) is a (positive) measure by 9.11, finite by 9.17. We prove (a) \(\Rightarrow\) (b) \(\Rightarrow\) (a) and (a) \(\Leftrightarrow\) (c).

Proof that (a) implies (b). Suppose \(\nu \ll \mu\). Let \(E \in \mathcal{S}\) with \(\mu(E) = 0\). We must show \(|\nu|(E) = 0\). Suppose \(n \in \mathbf{Z}^{+}\) and \(E_{1}, \ldots, E_{n}\) are disjoint sets in \(\mathcal{S}\) with \(E_{1} \cup \cdots \cup E_{n} \subseteq E\). For each \(k\) we have \(E_{k} \subseteq E\), and hence

\begin{equation*} \mu(E_{k}) \le \mu(E) = 0 \end{equation*}

by monotonicity of the (positive) measure \(\mu\). Thus \(\mu(E_{k}) = 0\), which by (a) gives \(\nu(E_{k}) = 0\). Therefore

\begin{equation*} |\nu(E_{1})| + \cdots + |\nu(E_{n})| = 0. \end{equation*}

Because this holds for every such finite disjoint family, the supremum defining \(|\nu|(E)\) in 9.8 equals \(0\). Hence \(|\nu|(E) = 0\), proving \(|\nu| \ll \mu\).

Proof that (b) implies (a). Suppose \(|\nu| \ll \mu\). Let \(E \in \mathcal{S}\) with \(\mu(E) = 0\). Then \(|\nu|(E) = 0\). Taking \(n = 1\) and \(E_{1} = E\) in the definition 9.8 of \(|\nu|\) shows that \(|\nu(E)| \le |\nu|(E)\) (this is the first of the bullet points following 9.8). Hence \(|\nu(E)| \le 0\), so \(\nu(E) = 0\). Thus \(\nu \ll \mu\).

Proof that (a) implies (c). Suppose \(\nu \ll \mu\) and \(E \in \mathcal{S}\) with \(\mu(E) = 0\). Then \(\nu(E) = 0\), and hence

\begin{equation*} (\operatorname{Re}\nu)(E) = \operatorname{Re}\bigl(\nu(E)\bigr) = 0 \qquad \text{and} \qquad (\operatorname{Im}\nu)(E) = \operatorname{Im}\bigl(\nu(E)\bigr) = 0. \end{equation*}

Thus \(\operatorname{Re}\nu \ll \mu\) and \(\operatorname{Im}\nu \ll \mu\).

Proof that (c) implies (a). Suppose \(\operatorname{Re}\nu \ll \mu\) and \(\operatorname{Im}\nu \ll \mu\), and suppose \(E \in \mathcal{S}\) with \(\mu(E) = 0\). Then \((\operatorname{Re}\nu)(E) = 0\) and \((\operatorname{Im}\nu)(E) = 0\), and hence

\begin{equation*} \nu(E) = (\operatorname{Re}\nu)(E) + i\,(\operatorname{Im}\nu)(E) = 0. \end{equation*}

Thus \(\nu \ll \mu\).

Having shown (a) \(\Leftrightarrow\) (b) and (a) \(\Leftrightarrow\) (c), all three statements are equivalent.

Problem (9B.11)

Suppose \(\mu\) is a (positive) measure on a measurable space \((X, \mathcal{S})\) and \(\nu\) is a real measure on \((X, \mathcal{S})\). Show that \(\nu \ll \mu\) if and only if \(\nu^{+} \ll \mu\) and \(\nu^{-} \ll \mu\).

Solution

Here \(\nu^{+}\) and \(\nu^{-}\) are the finite (positive) measures provided by the Jordan Decomposition Theorem (9.30), so that

\begin{equation*} \nu = \nu^{+} - \nu^{-}, \qquad \nu^{+} \perp \nu^{-}, \qquad |\nu| = \nu^{+} + \nu^{-}; \end{equation*}

the first two conditions are 9.31, and the third is the additional assertion of 9.30.

First suppose \(\nu^{+} \ll \mu\) and \(\nu^{-} \ll \mu\). Let \(E \in \mathcal{S}\) with \(\mu(E) = 0\). Then \(\nu^{+}(E) = 0\) and \(\nu^{-}(E) = 0\), and hence

\begin{equation*} \nu(E) = \nu^{+}(E) - \nu^{-}(E) = 0. \end{equation*}

Thus \(\nu \ll \mu\).

Conversely suppose \(\nu \ll \mu\). Let \(A, B \in \mathcal{S}\) be a Hahn decomposition of \(\nu\) as in 9.23, so that (as in the proof of 9.30)

\begin{equation*} \nu^{+}(E) = \nu(E \cap A) \qquad \text{and} \qquad \nu^{-}(E) = -\nu(E \cap B) \end{equation*}

for all \(E \in \mathcal{S}\). Now suppose \(E \in \mathcal{S}\) with \(\mu(E) = 0\). Because \(E \cap A \subseteq E\) and \(E \cap B \subseteq E\), monotonicity of the (positive) measure \(\mu\) gives

\begin{equation*} \mu(E \cap A) \le \mu(E) = 0 \qquad \text{and} \qquad \mu(E \cap B) \le \mu(E) = 0. \end{equation*}

Hence \(\mu(E \cap A) = \mu(E \cap B) = 0\), and so \(\nu \ll \mu\) implies \(\nu(E \cap A) = 0\) and \(\nu(E \cap B) = 0\). Therefore

\begin{equation*} \nu^{+}(E) = \nu(E \cap A) = 0 \qquad \text{and} \qquad \nu^{-}(E) = -\nu(E \cap B) = 0. \end{equation*}

Thus \(\nu^{+} \ll \mu\) and \(\nu^{-} \ll \mu\), completing the proof.

Alternatively, the forward implication follows from Exercise 10 of this section: \(\nu \ll \mu\) implies \(|\nu| \ll \mu\), and if \(\mu(E) = 0\) then \(\nu^{+}(E) + \nu^{-}(E) = |\nu|(E) = 0\) by 9.30, which forces \(\nu^{+}(E) = \nu^{-}(E) = 0\) because \(\nu^{+}\) and \(\nu^{-}\) take values in \([0, \infty)\).

Problem (9B.12)

Suppose \(\mu\) is a (positive) measure on a measurable space \((X, \mathcal{S})\). Prove that

\begin{equation*} \{\nu \in \mathcal{M}_{\mathbf{F}}(\mathcal{S}) : \nu \ll \mu\} \end{equation*}

is a closed subspace of \(\mathcal{M}_{\mathbf{F}}(\mathcal{S})\).

Solution

Write

\begin{equation*} V = \{\nu \in \mathcal{M}_{\mathbf{F}}(\mathcal{S}) : \nu \ll \mu\}, \end{equation*}

and let

\begin{equation*} \mathcal{N} = \{E \in \mathcal{S} : \mu(E) = 0\} \end{equation*}

be the collection of sets of \(\mu\)-measure \(0\). By definition 9.32,

\begin{equation*} V = \{\nu \in \mathcal{M}_{\mathbf{F}}(\mathcal{S}) : \nu(E) = 0 \text{ for every } E \in \mathcal{N}\}. \end{equation*}

Recall from 9.14 that \(\mathcal{M}_{\mathbf{F}}(\mathcal{S})\) is the vector space of real measures on \((X, \mathcal{S})\) if \(\mathbf{F} = \mathbf{R}\) and of complex measures if \(\mathbf{F} = \mathbf{C}\), with the operations of 9.13, namely

\begin{equation*} (\nu_{1} + \nu_{2})(E) = \nu_{1}(E) + \nu_{2}(E) \qquad \text{and} \qquad (\alpha\nu)(E) = \alpha\,\nu(E) \end{equation*}

for \(E \in \mathcal{S}\) and \(\alpha \in \mathbf{F}\); by 9.18 it is a Banach space with the total variation norm \(\|\nu\| = |\nu|(X)\) of 9.15.

First we show \(V\) is a subspace of \(\mathcal{M}_{\mathbf{F}}(\mathcal{S})\).

The zero measure lies in \(V\), because it assigns \(0\) to every set in \(\mathcal{S}\); in particular \(V \ne \varnothing\).

Suppose \(\nu_{1}, \nu_{2} \in V\) and \(E \in \mathcal{N}\). Then \(\nu_{1}(E) = \nu_{2}(E) = 0\), so

\begin{equation*} (\nu_{1} + \nu_{2})(E) = \nu_{1}(E) + \nu_{2}(E) = 0. \end{equation*}

As \(E \in \mathcal{N}\) was arbitrary, \(\nu_{1} + \nu_{2} \in V\). Thus \(V\) is closed under addition.

Suppose \(\nu \in V\), \(\alpha \in \mathbf{F}\), and \(E \in \mathcal{N}\). Then

\begin{equation*} (\alpha\nu)(E) = \alpha\,\nu(E) = \alpha \cdot 0 = 0. \end{equation*}

As \(E \in \mathcal{N}\) was arbitrary, \(\alpha\nu \in V\). Thus \(V\) is closed under scalar multiplication, and hence \(V\) is a subspace of \(\mathcal{M}_{\mathbf{F}}(\mathcal{S})\).

Now we show \(V\) is a closed subset of \(\mathcal{M}_{\mathbf{F}}(\mathcal{S})\) in the total variation norm.

Suppose \(\nu_{1}, \nu_{2}, \ldots\) is a sequence in \(V\) and \(\nu \in \mathcal{M}_{\mathbf{F}}(\mathcal{S})\) with

\begin{equation*} \lim_{k \to \infty} \|\nu - \nu_{k}\| = 0. \end{equation*}

Let \(E \in \mathcal{S}\) with \(\mu(E) = 0\). For each \(k \in \mathbf{Z}^{+}\) we have \(\nu_{k}(E) = 0\), and therefore

\begin{equation*} |\nu(E)| = |\nu(E) - \nu_{k}(E)| = \bigl|(\nu - \nu_{k})(E)\bigr| \le |\nu - \nu_{k}|(E) \le |\nu - \nu_{k}|(X) = \|\nu - \nu_{k}\|, \end{equation*}

where the first inequality is the bound \(|\lambda(E)| \le |\lambda|(E)\) noted after 9.8 (applied to the measure \(\lambda = \nu - \nu_{k}\)), and the second inequality holds because the total variation measure \(|\nu - \nu_{k}|\) is a (positive) measure by 9.11 and \(E \subseteq X\).

The right side of the display above tends to \(0\) as \(k \to \infty\), while the left side does not depend on \(k\). Hence \(|\nu(E)| = 0\), so \(\nu(E) = 0\). As \(E\) was an arbitrary set in \(\mathcal{S}\) with \(\mu(E) = 0\), we conclude that \(\nu \ll \mu\), that is, \(\nu \in V\).

Therefore \(V\) contains the limit of every convergent sequence of its elements, so \(V\) is closed. Hence \(V\) is a closed subspace of \(\mathcal{M}_{\mathbf{F}}(\mathcal{S})\), as desired. Because \(\mathcal{M}_{\mathbf{F}}(\mathcal{S})\) is a Banach space (9.18), it follows that \(V\) is itself a Banach space with the total variation norm.

Problem (9B.13)

Give an example to show that the Radon–Nikodym Theorem (9.36) can fail if the \(\sigma\)-finite hypothesis is eliminated.

Solution

Let \(X = [0, 1]\) and let \(\mathcal{S} = \mathcal{B}_{[0,1]}\) be the \(\sigma\)-algebra of Borel subsets of \([0, 1]\). Define

\begin{equation*} \mu(E) = \text{the number of elements of } E \end{equation*}

for \(E \in \mathcal{S}\), with \(\mu(E) = \infty\) when \(E\) is an infinite set; thus \(\mu\) is counting measure on \((X, \mathcal{S})\), which is a (positive) measure. Let \(\nu = \lambda|_{\mathcal{S}}\) be Lebesgue measure restricted to the Borel subsets of \([0, 1]\); because \(\nu([0,1]) = 1 < \infty\), the measure \(\nu\) is a finite (positive) measure and hence is a real measure, so in particular \(\nu\) is a complex measure on \((X, \mathcal{S})\).

We claim that \(\nu \ll \mu\) but that there is no \(h \in L^{1}(\mu)\) with \(d\nu = h \, d\mu\). Together with the observation that \(\mu\) is not \(\sigma\)-finite, this shows that the \(\sigma\)-finite hypothesis in 9.36 cannot be deleted.

First, \(\mu\) is not \(\sigma\)-finite. Indeed, if \(X = \bigcup_{k=1}^{\infty} X_{k}\) with each \(X_{k} \in \mathcal{S}\) and \(\mu(X_{k}) < \infty\), then each \(X_{k}\) is a finite set, so \(X\) would be a countable union of finite sets and hence countable. But \([0,1]\) is uncountable, a contradiction.

Second, \(\nu \ll \mu\). Suppose \(E \in \mathcal{S}\) and \(\mu(E) = 0\). By the definition of counting measure, \(\mu(E) = 0\) forces \(E = \varnothing\), and hence \(\nu(E) = \nu(\varnothing) = 0\). Thus every set of \(\mu\)-measure \(0\) has \(\nu\)-measure \(0\), which is exactly the statement \(\nu \ll \mu\) from 9.32. (The absolute continuity here is vacuous, but that is legitimate.)

Third, no \(h \in L^{1}(\mu)\) satisfies \(d\nu = h \, d\mu\). Suppose, for a contradiction, that \(h \in L^{1}(\mu)\) and

\begin{equation*} \nu(E) = \int_{E} h \, d\mu \qquad \text{for every } E \in \mathcal{S}. \end{equation*}

Fix \(x \in [0,1]\) and take \(E = \{x\}\), which is a Borel set. The function \(h\chi_{\{x\}}\) is identically equal to the simple function \(h(x)\chi_{\{x\}}\), because both vanish off \(\{x\}\) and both take the value \(h(x)\) at \(x\). Hence

\begin{equation*} \int_{\{x\}} h \, d\mu = \int h(x)\chi_{\{x\}} \, d\mu = h(x)\,\mu(\{x\}) = h(x), \end{equation*}

using \(\mu(\{x\}) = 1\). On the other hand \(\nu(\{x\}) = \lambda(\{x\}) = 0\). Hence \(h(x) = 0\). Because \(x \in [0,1]\) was arbitrary, \(h\) is identically \(0\) on \(X\). But then

\begin{equation*} 1 = \nu([0,1]) = \int_{[0,1]} h \, d\mu = 0, \end{equation*}

a contradiction.

Therefore \(\nu\) is a complex (indeed finite positive) measure with \(\nu \ll \mu\) for which no \(h \in L^{1}(\mu)\) gives \(d\nu = h \, d\mu\). The only hypothesis of 9.36 that fails is the \(\sigma\)-finiteness of \(\mu\), so that hypothesis cannot be eliminated.

The same example works with \(\mathcal{S}\) taken to be the \(\sigma\)-algebra of Lebesgue measurable subsets of \([0,1]\), or with \([0,1]\) replaced by any uncountable set carrying a nonzero finite measure that vanishes on singletons.

Problem (9B.14)

Suppose \(\mu\) is a (positive) \(\sigma\)-finite measure on a measurable space \((X, \mathcal{S})\) and \(\nu\) is a complex measure on \((X, \mathcal{S})\). Show that the following are equivalent.

(a)
\(\nu \ll \mu\).
(b)
For every \(\varepsilon > 0\), there exists \(\delta > 0\) such that \(|\nu(E)| < \varepsilon\) for every set \(E \in \mathcal{S}\) with \(\mu(E) < \delta\).
(c)
For every \(\varepsilon > 0\), there exists \(\delta > 0\) such that \(|\nu|(E) < \varepsilon\) for every set \(E \in \mathcal{S}\) with \(\mu(E) < \delta\).
Solution

We prove (c) \(\Rightarrow\) (b) \(\Rightarrow\) (a) \(\Rightarrow\) (c).

Proof that (c) implies (b). Suppose (c) holds and \(\varepsilon > 0\). Choose \(\delta > 0\) as in (c). If \(E \in \mathcal{S}\) and \(\mu(E) < \delta\), then

\begin{equation*} |\nu(E)| \le |\nu|(E) < \varepsilon, \end{equation*}

where the first inequality is the bound noted after 9.8 (take \(n = 1\) and \(E_{1} = E\) in the supremum defining \(|\nu|(E)\)). Thus (b) holds with the same \(\delta\).

Proof that (b) implies (a). Suppose (b) holds and \(E \in \mathcal{S}\) with \(\mu(E) = 0\). Let \(\varepsilon > 0\) and choose \(\delta > 0\) as in (b). Since \(\mu(E) = 0 < \delta\), we get \(|\nu(E)| < \varepsilon\). As \(\varepsilon > 0\) was arbitrary, \(\nu(E) = 0\). Hence \(\nu \ll \mu\), which is (a).

Proof that (a) implies (c). Suppose \(\nu \ll \mu\). By Exercise 10 of this section, \(|\nu| \ll \mu\). Recall also that \(|\nu|\) is a (positive) measure by 9.11 and that \(|\nu|(X) = \|\nu\| < \infty\) by 9.17; thus \(|\nu|\) is a finite (positive) measure and in particular is a real measure, so 9.7 applies to it.

Suppose (c) fails. Then there exists \(\varepsilon > 0\) such that for every \(\delta > 0\) there is a set \(E \in \mathcal{S}\) with \(\mu(E) < \delta\) but \(|\nu|(E) \ge \varepsilon\). Applying this with \(\delta = 2^{-k}\) for each \(k \in \mathbf{Z}^{+}\) produces sets \(E_{1}, E_{2}, \ldots \in \mathcal{S}\) with

\begin{equation*} \mu(E_{k}) < \frac{1}{2^{k}} \qquad \text{and} \qquad |\nu|(E_{k}) \ge \varepsilon . \end{equation*}

For \(n \in \mathbf{Z}^{+}\) let

\begin{equation*} F_{n} = \bigcup_{k = n}^{\infty} E_{k}, \end{equation*}

and let \(F = \bigcap_{n=1}^{\infty} F_{n}\). Note that \(F_{1} \supseteq F_{2} \supseteq \cdots\).

Countable subadditivity of the measure \(\mu\) (2.58) gives

\begin{equation*} \mu(F_{n}) \le \sum_{k = n}^{\infty} \mu(E_{k}) < \sum_{k = n}^{\infty} \frac{1}{2^{k}} = \frac{1}{2^{n-1}} \end{equation*}

for each \(n \in \mathbf{Z}^{+}\). Because \(F \subseteq F_{n}\), monotonicity of \(\mu\) gives \(\mu(F) \le 2^{-(n-1)}\) for every \(n\), and hence \(\mu(F) = 0\). Since \(|\nu| \ll \mu\), this implies

\begin{equation*} |\nu|(F) = 0 . \end{equation*}

On the other hand, \(E_{n} \subseteq F_{n}\) and \(|\nu|\) is a (positive) measure, so \(|\nu|(F_{n}) \ge |\nu|(E_{n}) \ge \varepsilon\) for every \(n \in \mathbf{Z}^{+}\). Because \(|\nu|\) is a finite (positive) measure and hence a real measure, 9.7(d) applied to the decreasing sequence \(F_{1} \supseteq F_{2} \supseteq \cdots\) gives

\begin{equation*} |\nu|(F) = \lim_{n \to \infty} |\nu|(F_{n}) \ge \varepsilon > 0 . \end{equation*}

This contradicts \(|\nu|(F) = 0\). Hence (c) holds, completing the cycle of implications and the proof.

Two remarks are worth recording.

First, the proof of (a) \(\Rightarrow\) (c) above never used \(\sigma\)-finiteness of \(\mu\); it used only that \(|\nu|\) is a finite measure, which holds for every complex measure by 9.17. So the three conditions are in fact equivalent for an arbitrary (positive) measure \(\mu\). The hypothesis of \(\sigma\)-finiteness is presumably included because it permits the following alternative, and perhaps more natural, argument for (a) \(\Rightarrow\) (c).

Second, here is that alternative argument. Assume \(\mu\) is \(\sigma\)-finite and \(\nu \ll \mu\). By the Radon–Nikodym Theorem (9.36) there exists \(h \in L^{1}(\mu)\) with \(d\nu = h \, d\mu\), and then 9.10 gives

\begin{equation*} |\nu|(E) = \int_{E} |h| \, d\mu \qquad \text{for every } E \in \mathcal{S}. \end{equation*}

Let \(\varepsilon > 0\). For \(n \in \mathbf{Z}^{+}\) set \(h_{n} = \min\{|h|, n\}\). Then \(h_{1} \le h_{2} \le \cdots\) and \(h_{n} \to |h|\) pointwise, so the Monotone Convergence Theorem (3.11) gives \(\int_{X} h_{n} \, d\mu \to \int_{X} |h| \, d\mu < \infty\). Choose \(n \in \mathbf{Z}^{+}\) with

\begin{equation*} \int_{X} \bigl(|h| - h_{n}\bigr) \, d\mu < \frac{\varepsilon}{2}, \end{equation*}

which is legitimate because \(|h| - h_{n} \ge 0\) and \(\int_{X}|h| \, d\mu\) is finite. Let \(\delta = \varepsilon/(2n)\). If \(E \in \mathcal{S}\) and \(\mu(E) < \delta\), then

\begin{equation*} |\nu|(E) = \int_{E} |h| \, d\mu = \int_{E} \bigl(|h| - h_{n}\bigr) \, d\mu + \int_{E} h_{n} \, d\mu \le \frac{\varepsilon}{2} + n\,\mu(E) < \frac{\varepsilon}{2} + \frac{\varepsilon}{2} = \varepsilon . \end{equation*}

Thus (c) holds.

Problem (9B.15)

Prove 9.42 [with the extra hypothesis that \(\mu\) is a \(\sigma\)-finite (positive) measure] in the case where \(p = 1\).

Solution

Throughout, \(\mu\) is a \(\sigma\)-finite (positive) measure on a measurable space \((X, \mathcal{S})\), and \(p = 1\), so that the dual exponent is \(p^{\prime} = \infty\). For \(h \in L^\infty(\mu)\) define \(\varphi_h : L^1(\mu) \to \mathbf{F}\) by

\begin{equation*} \varphi_h(f) = \int f h \, d\mu . \end{equation*}

We must show that \(h \mapsto \varphi_h\) is a one-to-one linear map from \(L^\infty(\mu)\) onto \(\bigl(L^1(\mu)\bigr)^{\prime}\) and that \(\|\varphi_h\| = \|h\|_\infty\) for every \(h \in L^\infty(\mu)\).

Step 1: each \(\varphi_h\) is a bounded linear functional, with \(\|\varphi_h\| \le \|h\|_\infty\).

If \(h \in L^\infty(\mu)\) and \(f \in L^1(\mu)\), then \(|f h| \le \|h\|_\infty |f|\) almost everywhere, so \(f h \in L^1(\mu)\) and

\begin{equation*} |\varphi_h(f)| \le \int |f h| \, d\mu \le \|h\|_\infty \|f\|_1 \end{equation*}

(this is 7.9 with the exponent pair \(1, \infty\)). Thus \(\varphi_h\) is a bounded linear functional on \(L^1(\mu)\) with \(\|\varphi_h\| \le \|h\|_\infty\). That \(h \mapsto \varphi_h\) is linear is immediate from the linearity of the integral in \(h\).

Step 2: \(\|\varphi_h\| = \|h\|_\infty\); consequently \(h \mapsto \varphi_h\) is one-to-one.

Let \(h \in L^\infty(\mu)\). If \(\|h\|_\infty = 0\) then \(h = 0\) in \(L^\infty(\mu)\) and both sides are \(0\), so assume \(\|h\|_\infty > 0\) and let \(c\) be any number with \(0 \le c < \|h\|_\infty\). Set

\begin{equation*} A = \{x \in X : |h(x)| > c\} . \end{equation*}

By the definition of the essential supremum, \(\mu(A) > 0\). Because \(\mu\) is \(\sigma\)-finite, there are sets \(X_1, X_2, \ldots \in \mathcal{S}\) with \(\bigcup_{k=1}^\infty X_k = X\) and \(\mu(X_k) < \infty\) for each \(k\). Since \(A = \bigcup_{k=1}^\infty (A \cap X_k)\) and \(\mu(A) > 0\), countable subadditivity gives some \(k\) with \(\mu(A \cap X_k) > 0\). Put \(B = A \cap X_k\), so

\begin{equation*} 0 < \mu(B) < \infty . \end{equation*}

Define \(f : X \to \mathbf{F}\) by \(f(x) = \overline{h(x)} / |h(x)|\) for \(x \in B\) (note \(|h(x)| > c \ge 0\) on \(B\), so this makes sense) and \(f(x) = 0\) for \(x \in X \setminus B\). Then \(|f| = \chi_B\), so \(f \in L^1(\mu)\) with \(\|f\|_1 = \mu(B) \in (0, \infty)\), and \(f h = \chi_B |h|\). Hence, using 3.8,

\begin{equation*} \varphi_h(f) = \int_B |h| \, d\mu \ge c \, \mu(B) = c \|f\|_1 . \end{equation*}

Therefore \(\|\varphi_h\| \ge c\). Letting \(c \uparrow \|h\|_\infty\) gives \(\|\varphi_h\| \ge \|h\|_\infty\), which combined with Step 1 yields \(\|\varphi_h\| = \|h\|_\infty\).

If now \(h_1, h_2 \in L^\infty(\mu)\) satisfy \(\varphi_{h_1} = \varphi_{h_2}\), then by linearity and the norm identity just proved,

\begin{equation*} \|h_1 - h_2\|_\infty = \|\varphi_{h_1 - h_2}\| = \|\varphi_{h_1} - \varphi_{h_2}\| = 0, \end{equation*}

so \(h_1 = h_2\) as elements of \(L^\infty(\mu)\). Thus \(h \mapsto \varphi_h\) is one-to-one.

Step 3: surjectivity when \(\mu\) is a finite measure.

Suppose \(\mu(X) < \infty\) and \(\varphi \in \bigl(L^1(\mu)\bigr)^{\prime}\). Because \(\mu\) is finite, \(\chi_E \in L^1(\mu)\) for every \(E \in \mathcal{S}\), so we may define \(\nu : \mathcal{S} \to \mathbf{F}\) by

\begin{equation*} \nu(E) = \varphi(\chi_E) . \end{equation*}

Clearly \(\nu(\emptyset) = \varphi(0) = 0\). Suppose \(E_1, E_2, \ldots\) are disjoint sets in \(\mathcal{S}\) with union \(E\). Then

\begin{equation*} \Bigl\| \chi_E - \sum_{k=1}^n \chi_{E_k} \Bigr\|_1 = \mu\Bigl( \bigcup_{k=n+1}^\infty E_k \Bigr) = \sum_{k=n+1}^\infty \mu(E_k) \longrightarrow 0 \end{equation*}

as \(n \to \infty\), because \(\sum_{k=1}^\infty \mu(E_k) = \mu(E) < \infty\) is a convergent series of nonnegative numbers. Since \(\varphi\) is a continuous linear functional,

\begin{equation*} \nu(E) = \varphi(\chi_E) = \lim_{n \to \infty} \varphi\Bigl( \sum_{k=1}^n \chi_{E_k} \Bigr) = \lim_{n \to \infty} \sum_{k=1}^n \nu(E_k) = \sum_{k=1}^\infty \nu(E_k) . \end{equation*}

Thus \(\nu\) is countably additive, so \(\nu\) is a complex measure on \((X, \mathcal{S})\) (and a real measure if \(\mathbf{F} = \mathbb{R}\)).

If \(E \in \mathcal{S}\) and \(\mu(E) = 0\), then \(\chi_E\) is the \(0\) element of \(L^1(\mu)\), so \(\nu(E) = \varphi(\chi_E) = 0\). Hence \(\nu \ll \mu\). Because \(\mu\) is finite (hence \(\sigma\)-finite), the Radon-Nikodym Theorem (9.36) provides \(h \in L^1(\mu)\) with \(d\nu = h \, d\mu\), that is,

\begin{equation*} \varphi(\chi_E) = \nu(E) = \int_E h \, d\mu = \int \chi_E h \, d\mu \end{equation*}

for every \(E \in \mathcal{S}\). By linearity of \(\varphi\) and of the integral,

\begin{equation*} \varphi(f) = \int f h \, d\mu \qquad \text{for every simple } \mathcal{S}\text{-measurable } f : X \to \mathbf{F} . \end{equation*}

Next we extend this identity to all bounded measurable \(f\). Let \(f \in L^\infty(\mu)\); after modifying \(f\) on a set of measure \(0\) we may assume \(f\) is bounded on \(X\). By 2.89 (applied to the real and imaginary parts of \(f\)) there is a sequence \(f_1, f_2, \ldots\) of simple \(\mathcal{S}\)-measurable functions converging to \(f\) uniformly on \(X\). Then \(\|f - f_k\|_1 \le \mu(X) \sup_X |f - f_k| \to 0\), so \(\varphi(f_k) \to \varphi(f)\) by continuity of \(\varphi\); and \(\bigl| \int (f - f_k) h \, d\mu \bigr| \le \|h\|_1 \sup_X |f - f_k| \to 0\) because \(h \in L^1(\mu)\). Hence

\begin{equation*} \varphi(f) = \int f h \, d\mu \qquad \text{for every } f \in L^\infty(\mu) . \end{equation*}

Now we show \(\|h\|_\infty \le \|\varphi\|\). Suppose not; then the set \(A = \{x \in X : |h(x)| > \|\varphi\|\}\) has \(\mu(A) > 0\). Writing \(A = \bigcup_{n=1}^\infty A_n\) with \(A_n = \{x \in X : |h(x)| \ge \|\varphi\| + \frac1n\}\), we get \(\mu(A_n) > 0\) for some \(n\); also \(\mu(A_n) \le \mu(X) < \infty\). Define \(f = \chi_{A_n} \overline{h} / |h|\) (with \(f = 0\) off \(A_n\); note \(|h| > 0\) on \(A_n\)). Then \(f\) is a bounded measurable function with \(|f| = \chi_{A_n}\), so \(\|f\|_1 = \mu(A_n)\), and \(f h = \chi_{A_n} |h|\). The displayed identity gives

\begin{equation*} \|\varphi\| \, \mu(A_n) = \|\varphi\| \, \|f\|_1 \ge |\varphi(f)| = \int_{A_n} |h| \, d\mu \ge \Bigl( \|\varphi\| + \tfrac1n \Bigr) \mu(A_n), \end{equation*}

which is impossible because \(0 < \mu(A_n) < \infty\). Hence \(h \in L^\infty(\mu)\) with \(\|h\|_\infty \le \|\varphi\|\).

Finally, let \(f \in L^1(\mu)\) be arbitrary and set \(f_k = f \chi_{\{x \in X : |f(x)| \le k\}}\). Then \(f_k \in L^\infty(\mu)\), \(|f - f_k| \le |f| \in L^1(\mu)\), and \(f_k(x) \to f(x)\) for almost every \(x\); so the Dominated Convergence Theorem (3.31) gives \(\|f - f_k\|_1 \to 0\). Since \(\varphi\) and \(\varphi_h\) are both continuous on \(L^1(\mu)\) (the latter by Step 1, as \(h \in L^\infty(\mu)\)) and they agree on each \(f_k\), we conclude \(\varphi(f) = \varphi_h(f)\). Thus \(\varphi = \varphi_h\), completing the proof when \(\mu\) is finite.

Step 4: surjectivity when \(\mu\) is \(\sigma\)-finite.

Let \(\varphi \in \bigl(L^1(\mu)\bigr)^{\prime}\). Choose \(X_1 \subseteq X_2 \subseteq \cdots\) in \(\mathcal{S}\) with \(\bigcup_{k=1}^\infty X_k = X\) and \(\mu(X_k) < \infty\) for each \(k\) (replace a given sequence by the increasing sequence of its partial unions). For \(E \in \mathcal{S}\) let \(\mathcal{S}_E = \{A \in \mathcal{S} : A \subseteq E\}\) and let \(\mu_E\) be the measure on \((E, \mathcal{S}_E)\) given by \(\mu_E(A) = \mu(A)\); as in the proof of 9.42, identify \(L^1(\mu_E)\) with the subspace of \(L^1(\mu)\) consisting of the functions vanishing almost everywhere outside \(E\), and let \(\varphi_E\) be the restriction of \(\varphi\) to that subspace, so \(\varphi_E\) is a bounded linear functional with \(\|\varphi_E\| \le \|\varphi\|\).

Because \(\mu_{X_k}\) is a finite measure, Step 3 gives \(h_k \in L^\infty(\mu_{X_k})\) with \(\|h_k\|_\infty \le \|\varphi\|\) and

\begin{equation*} \varphi(f) = \int_{X_k} f h_k \, d\mu \qquad \text{for all } f \in L^1(\mu_{X_k}) . \end{equation*}

If \(j < k\), then \(h_j\) and the restriction of \(h_k\) to \(X_j\) both represent \(\varphi_{X_j}\), so by the uniqueness (Step 2, applied on the measure space \((X_j, \mathcal{S}_{X_j}, \mu_{X_j})\)) we have \(h_j = h_k\) almost everywhere on \(X_j\). Hence there is an \(\mathcal{S}\)-measurable \(h : X \to \mathbf{F}\) such that for each \(k\) we have \(h = h_k\) almost everywhere on \(X_k\) (define \(h = h_1\) on \(X_1\) and \(h = h_k\) on \(X_k \setminus X_{k-1}\) for \(k \ge 2\)). Then \(|h| \le \|\varphi\|\) almost everywhere on each \(X_k\), hence almost everywhere on \(X\); thus \(h \in L^\infty(\mu)\) with \(\|h\|_\infty \le \|\varphi\|\).

Now let \(f \in L^1(\mu)\). Then \(|f - f \chi_{X_k}| \le |f|\) and \(f \chi_{X_k} \to f\) pointwise, so \(\|f - f \chi_{X_k}\|_1 \to 0\) by the Dominated Convergence Theorem (3.31). Also \(f \chi_{X_k} \in L^1(\mu_{X_k})\), so

\begin{equation*} \varphi(f) = \lim_{k \to \infty} \varphi(f \chi_{X_k}) = \lim_{k \to \infty} \int_{X_k} f h_k \, d\mu = \lim_{k \to \infty} \int f \chi_{X_k} h \, d\mu = \int f h \, d\mu, \end{equation*}

where the first equality uses the continuity of \(\varphi\), the third uses \(h = h_k\) almost everywhere on \(X_k\), and the last uses the Dominated Convergence Theorem again (with dominating function \(\|h\|_\infty |f| \in L^1(\mu)\)).

Thus \(\varphi = \varphi_h\), so \(h \mapsto \varphi_h\) maps \(L^\infty(\mu)\) onto \(\bigl(L^1(\mu)\bigr)^{\prime}\). Together with Steps 1 and 2 this proves 9.42 for \(p = 1\) under the hypothesis that \(\mu\) is \(\sigma\)-finite.

Problem (9B.16)

Explain where the proof of 9.42 fails if \(p = \infty\).

Solution

If \(p = \infty\), then \(p^{\prime} = 1\), so the assertion of 9.42 would be that every bounded linear functional on \(L^\infty(\mu)\) has the form \(f \mapsto \int f h \, d\mu\) for some \(h \in L^1(\mu)\).

The proof breaks down at its very first move: the construction of the measure \(\nu\). Given \(\varphi \in \bigl(L^\infty(\mu)\bigr)^{\prime}\) with \(\mu\) a finite measure, the proof sets \(\nu(E) = \varphi(\chi_E)\) and then verifies countable additivity by writing, for disjoint sets \(E_1, E_2, \ldots \in \mathcal{S}\) with union \(E\),

\begin{equation*} \varphi(\chi_E) = \varphi\Bigl( \sum_{k=1}^\infty \chi_{E_k} \Bigr) = \sum_{k=1}^\infty \varphi(\chi_{E_k}), \end{equation*}

where the crucial (and explicitly stated) point is that the series \(\sum_{k=1}^\infty \chi_{E_k}\) converges to \(\chi_E\) in the \(L^p(\mu)\)-norm, so that the continuity of \(\varphi\) may be applied. For \(1 \le p < \infty\) this convergence holds, because

\begin{equation*} \Bigl\| \chi_E - \sum_{k=1}^n \chi_{E_k} \Bigr\|_p = \Bigl\| \chi_{\bigcup_{k > n} E_k} \Bigr\|_p = \mu\Bigl( \bigcup_{k=n+1}^\infty E_k \Bigr)^{1/p} \longrightarrow 0 \end{equation*}

as \(n \to \infty\), since \(\sum_{k=1}^\infty \mu(E_k) = \mu(E) < \infty\). For \(p = \infty\) this computation collapses: the \(L^\infty(\mu)\)-norm of a characteristic function is \(1\) whenever the set has positive measure, so

\begin{equation*} \Bigl\| \chi_E - \sum_{k=1}^n \chi_{E_k} \Bigr\|_\infty = \Bigl\| \chi_{\bigcup_{k > n} E_k} \Bigr\|_\infty = 1 \end{equation*}

for every \(n\) such that \(\mu\bigl( \bigcup_{k > n} E_k \bigr) > 0\), which is the typical situation (for instance if every \(E_k\) has positive measure). Thus the partial sums do not converge to \(\chi_E\) in \(L^\infty(\mu)\), the continuity of \(\varphi\) gives no information, and all one can conclude is that \(E \mapsto \varphi(\chi_E)\) is finitely additive. A finitely additive set function is not a complex measure, so the Radon-Nikodym Theorem (9.36) cannot be applied, and with it every later step of the proof (the identity \(d\nu = h \, d\mu\), and hence 9.43 and 9.44) disappears.

A second, independent breakdown occurs in the part of the proof that reduces a general measure to sets of finite measure. There one chooses \(f_k \in L^p(\mu)\) with \(\|f_k\|_p \le 1\) and \(|\varphi(f_k)| > \|\varphi\| - \frac1k\), truncates \(f_k\) to \(f_k \chi_{\{x \in X : |f_k(x)| > 1/n\}}\) for large \(n\), and then asserts that \(\mu(D_k) < \infty\), where \(D_k = \{x \in X : |f_k(x)| > 1/n_k\}\), because \(f_k \in L^p(\mu)\). The truncation itself is harmless when \(p = \infty\): the difference \(f_k - f_k \chi_{\{|f_k| > 1/n\}}\) is bounded in absolute value by \(1/n\) at every point, so it tends to \(0\) even in the \(L^\infty(\mu)\)-norm. What is lost is the conclusion \(\mu(D_k) < \infty\). For \(p < \infty\) that conclusion comes from Markov’s inequality (4.1) applied to \(|f_k|^p \in L^1(\mu)\) with the constant \(n_k^{-p}\): since \(D_k \subseteq \{x \in X : |f_k(x)|^p \ge n_k^{-p}\}\), we get \(\mu(D_k) \le n_k^{\,p} \|f_k\|_p^p < \infty\). For \(p = \infty\) there is no such estimate, and membership in \(L^\infty(\mu)\) carries no information at all about the measure of the set where \(|f_k|\) is bounded away from \(0\) (indeed \(\|\chi_X\|_\infty = 1\) even when \(\mu(X) = \infty\)).

The same defect resurfaces in the final sentence of that proof, which asserts that every element of \(L^p(\mu)\) can be approximated in norm by elements of \(L^p(\mu_E)\) plus functions living on subsets of \(X \setminus E\) of finite measure; all such approximants vanish outside a \(\sigma\)-finite set. For \(p < \infty\) this is fine, but for \(p = \infty\) it fails: if \(\mu\) is counting measure on an uncountable set \(X\), then a set of finite measure is finite and hence a \(\sigma\)-finite subset of \(X\) is countable. So if \(g\) vanishes outside a countable set \(F\), then \(|\chi_X - g| = 1\) at every point of the uncountable set \(X \setminus F\), each of whose points has measure \(1\); hence \(\|\chi_X - g\|_\infty \ge 1\) for every approximant \(g\) of the permitted kind, and \(\chi_X \in L^\infty(\mu)\) is not approximated at all.

These are not merely deficiencies of this particular argument: the conclusion itself is false for \(p = \infty\). Let \(\lambda\) be Lebesgue measure on the Borel subsets of \([0,1]\), and regard \(C([0,1])\) as a subspace of \(L^\infty(\lambda)\). This is legitimate: distinct continuous functions on \([0,1]\) are not equal almost everywhere, since if \(f\) is continuous and \(|f(x_0)| > c\) then \(|f| > c\) on a nonempty relatively open interval about \(x_0\), a set of positive measure; the same observation shows that for continuous \(f\) the essential supremum of \(|f|\) equals \(\sup_{[0,1]} |f|\), so the inclusion of \(C([0,1])\) into \(L^\infty(\lambda)\) preserves norms.

Define \(\psi : C([0,1]) \to \mathbf{F}\) by \(\psi(f) = f(0)\). Then \(\psi\) is linear and \(|\psi(f)| \le \sup_{[0,1]} |f| = \|f\|_\infty\), so \(\psi\) is a bounded linear functional on this subspace of \(L^\infty(\lambda)\). By the Hahn-Banach Theorem (6.69), \(\psi\) extends to a bounded linear functional \(\varphi\) on \(L^\infty(\lambda)\).

Suppose there existed \(h \in L^1(\lambda)\) with \(\varphi(f) = \int f h \, d\lambda\) for all \(f \in L^\infty(\lambda)\). For \(n \in \mathbb{Z}^+\) let

\begin{equation*} g_n(x) = \max\{0, \, 1 - nx\} \qquad (0 \le x \le 1). \end{equation*}

Each \(g_n\) is continuous with \(g_n(0) = 1\), so \(\varphi(g_n) = \psi(g_n) = 1\) for every \(n\). On the other hand \(|g_n h| \le |h| \in L^1(\lambda)\) and \(g_n(x) \to 0\) for every \(x \in (0,1]\), so the Dominated Convergence Theorem (3.31) gives

\begin{equation*} 1 = \varphi(g_n) = \int g_n h \, d\lambda \longrightarrow 0, \end{equation*}

a contradiction. Hence no such \(h\) exists, and the map \(h \mapsto \varphi_h\) from \(L^1(\lambda)\) into \(\bigl(L^\infty(\lambda)\bigr)^{\prime}\) is not surjective. So 9.42 genuinely fails for \(p = \infty\), and no repair of the proof is possible.

Problem (9B.17)

Prove that if \(\mu\) is a (positive) measure and \(1 < p < \infty\), then \(L^p(\mu)\) is reflexive. [See the definition before Exercise 19 in Section 7B for the meaning of reflexive.]

Solution

Recall the definitions. If \(V\) is a normed vector space, then \(V^{\prime\prime} = (V^{\prime})^{\prime}\) is its double dual, and the canonical isometry \(\Phi : V \to V^{\prime\prime}\) is defined by

\begin{equation*} (\Phi f)(\varphi) = \varphi(f) \qquad \text{for } f \in V \text{ and } \varphi \in V^{\prime} ; \end{equation*}

\(\Phi\) is linear and \(\|\Phi f\| = \|f\|\) for every \(f \in V\) (Exercise 20 in Section 6D). A Banach space is reflexive if its canonical isometry is surjective. By 7.24, \(L^p(\mu)\) is a Banach space, so the question makes sense; we must show that \(\Phi\) maps \(L^p(\mu)\) onto \(L^p(\mu)^{\prime\prime}\).

Suppose \(\mu\) is a (positive) measure on \((X, \mathcal{S})\) and \(1 < p < \infty\). Let \(p^{\prime}\) be the dual exponent, so \(\frac1p + \frac1{p^{\prime}} = 1\). Because \(1 < p < \infty\) we also have \(1 < p^{\prime} < \infty\), and the dual exponent of \(p^{\prime}\) is \(p\) (both statements follow at once from \(\frac1{p^{\prime}} = 1 - \frac1p \in (0,1)\)).

Since \(1 < p < \infty\) and \(1 < p^{\prime} < \infty\), 9.42 applies with exponent \(p\) and also with exponent \(p^{\prime}\) (no \(\sigma\)-finiteness hypothesis is needed, because that hypothesis is required only when the exponent equals \(1\)). Thus we have two maps

\begin{equation*} \Lambda : L^{p^{\prime}}(\mu) \to \bigl(L^p(\mu)\bigr)^{\prime}, \qquad (\Lambda h)(f) = \int f h \, d\mu \quad \text{for } f \in L^p(\mu), \end{equation*}

\begin{equation*} \Gamma : L^{p}(\mu) \to \bigl(L^{p^{\prime}}(\mu)\bigr)^{\prime}, \qquad (\Gamma g)(h) = \int h g \, d\mu \quad \text{for } h \in L^{p^{\prime}}(\mu), \end{equation*}

and 9.42 tells us that each of \(\Lambda\) and \(\Gamma\) is a one-to-one linear map onto the indicated dual space, with \(\|\Lambda h\| = \|h\|_{p^{\prime}}\) and \(\|\Gamma g\| = \|g\|_p\). In particular \(\Lambda\) is a linear bijection with \(\|\Lambda h\| = \|h\|_{p^{\prime}}\) for all \(h\).

Now let \(\Psi \in L^p(\mu)^{\prime\prime} = \bigl( (L^p(\mu))^{\prime} \bigr)^{\prime}\) be arbitrary. Consider the composition

\begin{equation*} \Psi \circ \Lambda : L^{p^{\prime}}(\mu) \to \mathbf{F} . \end{equation*}

It is linear, being a composition of linear maps, and it is bounded, since

\begin{equation*} |(\Psi \circ \Lambda)(h)| \le \|\Psi\| \, \|\Lambda h\| = \|\Psi\| \, \|h\|_{p^{\prime}} \end{equation*}

for every \(h \in L^{p^{\prime}}(\mu)\). Hence \(\Psi \circ \Lambda \in \bigl(L^{p^{\prime}}(\mu)\bigr)^{\prime}\). Because \(\Gamma\) is onto \(\bigl(L^{p^{\prime}}(\mu)\bigr)^{\prime}\), there exists \(g \in L^p(\mu)\) with \(\Gamma g = \Psi \circ \Lambda\), that is,

\begin{equation*} \Psi(\Lambda h) = \int h g \, d\mu \qquad \text{for every } h \in L^{p^{\prime}}(\mu) . \end{equation*}

We claim \(\Phi g = \Psi\). Let \(\varphi \in \bigl(L^p(\mu)\bigr)^{\prime}\). Because \(\Lambda\) is onto \(\bigl(L^p(\mu)\bigr)^{\prime}\), there exists \(h \in L^{p^{\prime}}(\mu)\) with \(\varphi = \Lambda h\). Then

\begin{equation*} (\Phi g)(\varphi) = \varphi(g) = (\Lambda h)(g) = \int g h \, d\mu = \int h g \, d\mu = \Psi(\Lambda h) = \Psi(\varphi), \end{equation*}

where the third equality is the definition of \(\Lambda\) (evaluated at \(g \in L^p(\mu)\)) and the fifth is the displayed identity above. Since \(\varphi \in \bigl(L^p(\mu)\bigr)^{\prime}\) was arbitrary, \(\Phi g = \Psi\).

Thus every element of \(L^p(\mu)^{\prime\prime}\) lies in the range of the canonical isometry \(\Phi\), so \(\Phi\) is surjective and \(L^p(\mu)\) is reflexive.

Two remarks on the hypotheses. First, the argument uses the surjectivity assertion of 9.42 twice, once for each of the exponents \(p\) and \(p^{\prime}\); this is exactly why the hypothesis \(1 < p < \infty\) is needed, since it is what guarantees that the conjugate exponent \(p^{\prime}\) also lies strictly between \(1\) and \(\infty\). Second, note that only the linearity and surjectivity of \(\Lambda\) and \(\Gamma\), together with the boundedness of \(\Lambda\) (which follows from the norm identity \(\|\Lambda h\| = \|h\|_{p^{\prime}}\), and is what makes \(\Psi \circ \Lambda\) a bounded functional), were used; the injectivity of these maps was not needed.

Problem (9B.18)

Prove that \(L^1(\mathbb{R})\) is not reflexive.

Solution

Here \(L^1(\mathbb{R}) = L^1(\lambda)\), where \(\lambda\) denotes Lebesgue measure on the \(\sigma\)-algebra \(\mathcal{B}\) of Borel subsets of \(\mathbb{R}\) (by 3.45 one may equally take Lebesgue measure on the Lebesgue measurable subsets of \(\mathbb{R}\); nothing below changes); by 7.24 this is a Banach space. Write \(\Phi : L^1(\mathbb{R}) \to L^1(\mathbb{R})^{\prime\prime}\) for the canonical isometry, \((\Phi f)(\varphi) = \varphi(f)\). We must produce an element of \(L^1(\mathbb{R})^{\prime\prime}\) that is not in the range of \(\Phi\).

Because \(\lambda\) is a \(\sigma\)-finite measure (for instance \(\mathbb{R} = \bigcup_{k=1}^\infty (-k, k)\) with \(\lambda\bigl( (-k,k) \bigr) = 2k < \infty\)), Exercise 15 above (that is, 9.42 in the case \(p = 1\)) applies: the map

\begin{equation*} \Lambda : L^\infty(\mathbb{R}) \to \bigl(L^1(\mathbb{R})\bigr)^{\prime}, \qquad (\Lambda h)(f) = \int f h \, d\lambda , \end{equation*}

is a one-to-one linear map from \(L^\infty(\mathbb{R})\) onto \(\bigl(L^1(\mathbb{R})\bigr)^{\prime}\) with \(\|\Lambda h\| = \|h\|_\infty\) for every \(h \in L^\infty(\mathbb{R})\). Hence \(\Lambda^{-1} : \bigl(L^1(\mathbb{R})\bigr)^{\prime} \to L^\infty(\mathbb{R})\) is a linear bijection that also preserves norms.

Let \(C_b(\mathbb{R})\) denote the set of bounded continuous functions from \(\mathbb{R}\) to \(\mathbf{F}\), regarded as a subset of \(L^\infty(\mathbb{R})\). This is legitimate, and \(C_b(\mathbb{R})\) is a subspace of \(L^\infty(\mathbb{R})\) on which the \(L^\infty\)-norm is the supremum norm. Indeed, if \(h\) is continuous and \(|h(x_0)| > c\) for some \(x_0 \in \mathbb{R}\), then by continuity \(|h| > c\) on some open interval containing \(x_0\), and every nonempty open interval has positive Lebesgue measure; hence

\begin{equation*} \|h\|_\infty = \operatorname{ess\,sup} |h| = \sup_{\mathbb{R}} |h| \qquad \text{for } h \in C_b(\mathbb{R}) . \end{equation*}

In particular, two continuous functions that agree almost everywhere agree everywhere, so a continuous function is determined by its equivalence class in \(L^\infty(\mathbb{R})\), and evaluation at a point is a well-defined function on \(C_b(\mathbb{R}) \subseteq L^\infty(\mathbb{R})\).

Define \(\psi_0 : C_b(\mathbb{R}) \to \mathbf{F}\) by \(\psi_0(h) = h(0)\). Then \(\psi_0\) is linear and \(|\psi_0(h)| \le \sup_{\mathbb{R}} |h| = \|h\|_\infty\), so \(\psi_0\) is a bounded linear functional on the subspace \(C_b(\mathbb{R})\) of the normed vector space \(L^\infty(\mathbb{R})\). By the Hahn-Banach Theorem (6.69), \(\psi_0\) extends to a bounded linear functional \(\psi\) on all of \(L^\infty(\mathbb{R})\), with \(\|\psi\| = \|\psi_0\| = 1\).

Now define

\begin{equation*} \Psi = \psi \circ \Lambda^{-1} : \bigl(L^1(\mathbb{R})\bigr)^{\prime} \to \mathbf{F} . \end{equation*}

Then \(\Psi\) is linear, and for every \(\varphi \in \bigl(L^1(\mathbb{R})\bigr)^{\prime}\) we have \(|\Psi(\varphi)| \le \|\psi\| \, \|\Lambda^{-1}\varphi\|_\infty = \|\psi\| \, \|\varphi\|\), so \(\Psi\) is bounded. Thus \(\Psi \in L^1(\mathbb{R})^{\prime\prime}\).

Suppose, for the sake of contradiction, that \(L^1(\mathbb{R})\) is reflexive. Then there exists \(f \in L^1(\mathbb{R})\) with \(\Phi f = \Psi\). Let \(h \in L^\infty(\mathbb{R})\) be arbitrary and take \(\varphi = \Lambda h\). Then

\begin{equation*} \int f h \, d\lambda = (\Lambda h)(f) = (\Phi f)(\Lambda h) = \Psi(\Lambda h) = \psi\bigl( \Lambda^{-1}(\Lambda h) \bigr) = \psi(h) . \end{equation*}

So

\begin{equation*} \int f h \, d\lambda = \psi(h) \qquad \text{for every } h \in L^\infty(\mathbb{R}) . \end{equation*}

Now for \(n \in \mathbb{Z}^+\) define \(h_n : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} h_n(x) = \max\{ 0, \, 1 - n|x| \} . \end{equation*}

Each \(h_n\) is continuous and bounded (with \(0 \le h_n \le 1\)), so \(h_n \in C_b(\mathbb{R}) \subseteq L^\infty(\mathbb{R})\), and \(h_n(0) = 1\). Because \(\psi\) extends \(\psi_0\),

\begin{equation*} \psi(h_n) = \psi_0(h_n) = h_n(0) = 1 \qquad \text{for every } n \in \mathbb{Z}^+ . \end{equation*}

On the other hand, \(|f h_n| \le |f| \in L^1(\mathbb{R})\) for every \(n\), and \(h_n(x) \to 0\) for every \(x \ne 0\), so \(f(x) h_n(x) \to 0\) for almost every \(x \in \mathbb{R}\). The Dominated Convergence Theorem (3.31) therefore gives

\begin{equation*} \lim_{n \to \infty} \int f h_n \, d\lambda = 0 . \end{equation*}

Combining the last two displays with the identity \(\int f h \, d\lambda = \psi(h)\) applied to \(h = h_n\) yields

\begin{equation*} 1 = \lim_{n \to \infty} \psi(h_n) = \lim_{n \to \infty} \int f h_n \, d\lambda = 0, \end{equation*}

a contradiction.

Hence no \(f \in L^1(\mathbb{R})\) satisfies \(\Phi f = \Psi\); the canonical isometry \(\Phi\) is not surjective, and therefore \(L^1(\mathbb{R})\) is not reflexive.

The underlying reason is worth stating: composing with \(\Lambda^{-1}\) identifies \(\bigl(L^\infty(\mathbb{R})\bigr)^{\prime}\) with \(L^1(\mathbb{R})^{\prime\prime}\), and under that identification \(\psi\) is an element of the double dual that behaves like evaluation at \(0\), something no integration against an \(L^1\) function can imitate, since a single point has Lebesgue measure \(0\). This is the same phenomenon that makes 9.42 fail for \(p = \infty\) (Exercise 16).

Linear Maps on Hilbert Spaces

Exercises 10A

Problem (10A.1)

Define \(T : \ell^2 \to \ell^2\) by \(T(a_1, a_2, \ldots) = (0, a_1, a_2, \ldots)\). Find a formula for \(T^*\).

Solution

The answer is the left shift:

\begin{equation*} T^*(b_1, b_2, b_3, \ldots) = (b_2, b_3, b_4, \ldots). \end{equation*}

To verify this, first note that \(T\) is a bounded linear map (indeed an isometry, since \(\|Ta\| = \|a\|\) for every \(a \in \ell^2\)), so its adjoint \(T^* : \ell^2 \to \ell^2\) is defined by 10.1.

Let \(S : \ell^2 \to \ell^2\) be defined by \(S(b_1, b_2, b_3, \ldots) = (b_2, b_3, b_4, \ldots)\). Then \(S\) is a bounded linear map, because

\begin{equation*} \|Sb\|^2 = \sum_{n=2}^{\infty} |b_n|^2 \le \sum_{n=1}^{\infty} |b_n|^2 = \|b\|^2 . \end{equation*}

Now suppose \(a = (a_1, a_2, \ldots) \in \ell^2\) and \(b = (b_1, b_2, \ldots) \in \ell^2\). The \(n^{\text{th}}\) coordinate of \(Ta\) equals \(0\) if \(n = 1\) and equals \(a_{n-1}\) if \(n \ge 2\). Hence

\begin{equation*} \begin{aligned} \langle Ta, b \rangle &= \sum_{n=1}^{\infty} (Ta)_n \overline{b_n} = \sum_{n=2}^{\infty} a_{n-1} \overline{b_n} = \sum_{k=1}^{\infty} a_k \overline{b_{k+1}} \\ &= \sum_{k=1}^{\infty} a_k \overline{(Sb)_k} = \langle a, Sb \rangle . \end{aligned} \end{equation*}

All the series above converge absolutely by the Cauchy–Schwarz inequality (8.11), so the rearrangement of indices is legitimate.

Because \(\langle Ta, b \rangle = \langle a, Sb \rangle\) for all \(a, b \in \ell^2\), and because the element \(T^*b\) is by definition the unique element of \(\ell^2\) satisfying \(\langle Ta, b \rangle = \langle a, T^*b \rangle\) for all \(a \in \ell^2\) (see the discussion following 10.1, which uses the Riesz Representation Theorem 8.47), we conclude that \(T^* = S\). In other words,

\begin{equation*} T^*(b_1, b_2, b_3, \ldots) = (b_2, b_3, b_4, \ldots), \end{equation*}

so the adjoint of the right shift \(T\) is the left shift \(S\); these are the two operators of Example 10.27.

Problem (10A.2)

Suppose \(V\) is a Hilbert space, \(U\) is a closed subspace of \(V\), and \(T : U \to V\) is defined by \(Tf = f\). Describe the linear operator \(T^* : V \to U\).

Solution

The answer: \(T^*\) is the orthogonal projection \(P_U\) of \(V\) onto \(U\). In other words,

\begin{equation*} T^*g = P_U g \qquad \text{for every } g \in V . \end{equation*}

First we check that the setup makes sense. Because \(U\) is a closed subspace of the Hilbert space \(V\), it is complete (a closed subset of a complete metric space is complete), and thus \(U\) is itself a Hilbert space with the inner product it inherits from \(V\). The inclusion map \(T : U \to V\) is clearly linear, and \(\|Tf\| = \|f\|\) for every \(f \in U\); hence \(T\) is bounded and \(T \in \mathcal{B}(U, V)\). Therefore the adjoint \(T^* : V \to U\) is defined by 10.1.

Now suppose \(g \in V\). By 8.43, we can write

\begin{equation*} g = P_U g + (g - P_U g), \qquad P_U g \in U, \quad g - P_U g \in U^{\perp} , \end{equation*}

and \(P_U g \in U\) by the definition of the orthogonal projection. If \(f \in U\), then \(\langle f, g - P_U g \rangle = \overline{\langle g - P_U g, f \rangle} = 0\) because \(g - P_U g\) is orthogonal to every element of \(U\) [8.37(a)]. Hence

\begin{equation*} \langle Tf, g \rangle = \langle f, g \rangle = \langle f, P_U g \rangle + \langle f, g - P_U g \rangle = \langle f, P_U g \rangle . \end{equation*}

Thus \(\langle Tf, g \rangle = \langle f, P_U g \rangle\) for every \(f \in U\), with \(P_U g \in U\). Because \(T^*g\) is by definition the unique element of \(U\) satisfying \(\langle Tf, g \rangle = \langle f, T^*g \rangle\) for all \(f \in U\), we conclude that

\begin{equation*} T^* g = P_U g . \end{equation*}

So \(T^*\) is the orthogonal projection onto \(U\), regarded as a map from \(V\) onto \(U\) (it is onto because \(\operatorname{range} P_U = U\) by 8.45(a)). Note also that \(T^*T = I\) on \(U\), since \(P_U f = f\) for \(f \in U\); consistent with 10.13, we have \(\operatorname{null} T^* = U^{\perp} = (\operatorname{range} T)^{\perp}\).

Problem (10A.3)

Suppose \(V\) and \(W\) are Hilbert spaces and \(g \in V\), \(h \in W\). Define \(T \in \mathcal{B}(V, W)\) by \(Tf = \langle f, g \rangle h\). Find a formula for \(T^*\).

Solution

The answer:

\begin{equation*} T^* u = \langle u, h \rangle g \qquad \text{for every } u \in W . \end{equation*}

Thus \(T^*\) has the same form as \(T\), with the roles of \(g\) and \(h\) interchanged.

First note that \(T\) really is a bounded linear map. Linearity of \(f \mapsto \langle f, g \rangle h\) follows from linearity of the inner product in its first slot, and the Cauchy–Schwarz inequality (8.11) gives

\begin{equation*} \|Tf\| = |\langle f, g \rangle| \, \|h\| \le \|g\| \, \|h\| \, \|f\| \end{equation*}

for every \(f \in V\); hence \(T \in \mathcal{B}(V, W)\) with \(\|T\| \le \|g\| \, \|h\|\). Therefore \(T^* : W \to V\) is defined by 10.1.

Now define \(S : W \to V\) by \(Su = \langle u, h \rangle g\). By the same reasoning, \(S\) is a bounded linear map. Suppose \(f \in V\) and \(u \in W\). Then

\begin{equation*} \begin{aligned} \langle Tf, u \rangle &= \langle \langle f, g \rangle h, u \rangle = \langle f, g \rangle \langle h, u \rangle \\ &= \overline{\overline{\langle h, u \rangle}} \, \langle f, g \rangle = \overline{\langle u, h \rangle} \, \langle f, g \rangle \\ &= \langle f, \langle u, h \rangle g \rangle = \langle f, Su \rangle , \end{aligned} \end{equation*}

where we used homogeneity in the first slot, conjugate symmetry of the inner product, and conjugate homogeneity in the second slot.

Because \(T^*u\) is the unique element of \(V\) with \(\langle Tf, u \rangle = \langle f, T^*u \rangle\) for all \(f \in V\), we conclude that \(T^* = S\); that is,

\begin{equation*} T^* u = \langle u, h \rangle g \qquad \text{for all } u \in W . \end{equation*}

As a sanity check with 10.13: if \(g \ne 0\) and \(h \ne 0\), then \(\operatorname{range} T = \{ \alpha h : \alpha \in \mathbf{F} \}\), and indeed \(\operatorname{null} T^* = \{ u \in W : \langle u, h \rangle = 0 \} = \{h\}^{\perp} = (\operatorname{range} T)^{\perp}\), as 10.13(a) requires. Also \(\|T\| = \|T^*\| = \|g\| \, \|h\|\), in agreement with 10.11 (take \(f = g\) to see that the bound \(\|g\| \, \|h\|\) above is attained).

Problem (10A.4)

Suppose \(V\) and \(W\) are Hilbert spaces and \(T \in \mathcal{B}(V, W)\) has finite-dimensional range. Prove that \(T^*\) also has finite-dimensional range.

Solution

Let \(n = \dim \operatorname{range} T < \infty\). If \(n = 0\), then \(T = 0\), hence \(T^* = 0\) and \(\operatorname{range} T^* = \{0\}\) is finite-dimensional. So assume \(n \ge 1\).

The subspace \(\operatorname{range} T\) of \(W\) is an \(n\)-dimensional inner product space with the inner product inherited from \(W\). Applying the Gram–Schmidt procedure to any basis of \(\operatorname{range} T\) produces an orthonormal basis \(g_1, \ldots, g_n\) of \(\operatorname{range} T\). Thus

\begin{equation*} v = \sum_{j=1}^{n} \langle v, g_j \rangle g_j \qquad \text{for every } v \in \operatorname{range} T . \end{equation*}

Now suppose \(u \in W\). We claim that

\begin{equation*} T^* u = \sum_{j=1}^{n} \langle u, g_j \rangle T^* g_j . \end{equation*}

To prove this, let \(f \in V\) be arbitrary. Because \(Tf \in \operatorname{range} T\), the expansion above gives \(Tf = \sum_{j=1}^{n} \langle Tf, g_j \rangle g_j\). Hence

\begin{equation*} \begin{aligned} \langle f, T^* u \rangle &= \langle Tf, u \rangle = \Big\langle \sum_{j=1}^{n} \langle Tf, g_j \rangle g_j, \, u \Big\rangle = \sum_{j=1}^{n} \langle Tf, g_j \rangle \langle g_j, u \rangle \\ &= \sum_{j=1}^{n} \langle f, T^* g_j \rangle \, \overline{\langle u, g_j \rangle} = \Big\langle f, \, \sum_{j=1}^{n} \langle u, g_j \rangle T^* g_j \Big\rangle , \end{aligned} \end{equation*}

where in the fourth equality we used \(\langle Tf, g_j \rangle = \langle f, T^* g_j \rangle\) (the defining property 10.2 of the adjoint) and \(\langle g_j, u \rangle = \overline{\langle u, g_j \rangle}\).

Thus \(\big\langle f, \, T^* u - \sum_{j=1}^{n} \langle u, g_j \rangle T^* g_j \big\rangle = 0\) for every \(f \in V\). Taking \(f\) to be the element inside the second slot shows that this element has norm \(0\), and hence

\begin{equation*} T^* u = \sum_{j=1}^{n} \langle u, g_j \rangle T^* g_j , \end{equation*}

proving the claim.

The displayed formula shows that \(T^* u \in \operatorname{span}(T^* g_1, \ldots, T^* g_n)\) for every \(u \in W\). Therefore

\begin{equation*} \operatorname{range} T^* \subseteq \operatorname{span}(T^* g_1, \ldots, T^* g_n), \end{equation*}

and consequently \(\dim \operatorname{range} T^* \le n < \infty\). Thus \(T^*\) has finite-dimensional range.

Remark. Applying this result to \(T^*\) in place of \(T\), and using \((T^*)^* = T\) (see 10.11), gives \(\dim \operatorname{range} T \le \dim \operatorname{range} T^*\) as well. Hence \(\dim \operatorname{range} T^* = \dim \operatorname{range} T\) whenever either is finite.

Problem (10A.5)

Prove or give a counterexample: If \(V\) is a Hilbert space and \(T : V \to V\) is a bounded linear map such that \(\dim \operatorname{null} T < \infty\), then \(\dim \operatorname{null} T^* < \infty\).

Solution

The statement is false. Here is a counterexample.

Let \(V = \ell^2\) and define \(T : \ell^2 \to \ell^2\) by

\begin{equation*} T(a_1, a_2, a_3, \ldots) = (0, a_1, 0, a_2, 0, a_3, \ldots) ; \end{equation*}

in other words, \((Ta)_{2k} = a_k\) and \((Ta)_{2k-1} = 0\) for every \(k \in \mathbf{Z}^+\). In words: \(T\) spreads the entries of \(a\) over the even coordinates and puts \(0\) in every odd coordinate.

\(T\) is a bounded operator on \(\ell^2\). Linearity is clear from the coordinatewise formula. Furthermore, for every \(a \in \ell^2\),

\begin{equation*} \|Ta\|^2 = \sum_{k=1}^{\infty} |(Ta)_{2k}|^2 + \sum_{k=1}^{\infty} |(Ta)_{2k-1}|^2 = \sum_{k=1}^{\infty} |a_k|^2 = \|a\|^2 , \end{equation*}

so \(Ta \in \ell^2\) and \(T\) is an isometry; in particular \(T \in \mathcal{B}(\ell^2)\) with \(\|T\| = 1\).

\(\dim \operatorname{null} T = 0 < \infty\). Indeed, if \(Ta = 0\) then \(\|a\| = \|Ta\| = 0\), so \(a = 0\). Thus \(\operatorname{null} T = \{0\}\).

Computation of \(T^*\). Define \(S : \ell^2 \to \ell^2\) by \(Sb = (b_2, b_4, b_6, \ldots)\), i.e. \((Sb)_k = b_{2k}\). Then \(S\) is linear and \(\|Sb\|^2 = \sum_{k=1}^{\infty} |b_{2k}|^2 \le \|b\|^2\), so \(S\) is a bounded operator. For all \(a, b \in \ell^2\),

\begin{equation*} \langle Ta, b \rangle = \sum_{n=1}^{\infty} (Ta)_n \overline{b_n} = \sum_{k=1}^{\infty} a_k \overline{b_{2k}} = \sum_{k=1}^{\infty} a_k \overline{(Sb)_k} = \langle a, Sb \rangle , \end{equation*}

where the odd-indexed terms drop out because \((Ta)_{2k-1} = 0\), and all series converge absolutely by the Cauchy–Schwarz inequality (8.11). Since \(T^*b\) is the unique element of \(\ell^2\) with \(\langle Ta, b \rangle = \langle a, T^*b \rangle\) for all \(a \in \ell^2\), we get \(T^* = S\):

\begin{equation*} T^*(b_1, b_2, b_3, \ldots) = (b_2, b_4, b_6, \ldots) . \end{equation*}

\(\dim \operatorname{null} T^* = \infty\). From the formula for \(T^*\),

\begin{equation*} \operatorname{null} T^* = \{ b \in \ell^2 : b_{2k} = 0 \text{ for every } k \in \mathbf{Z}^+ \} , \end{equation*}

which is the set of elements of \(\ell^2\) whose even coordinates all vanish. For \(n \in \mathbf{Z}^+\) let \(e_n \in \ell^2\) be the vector with \(1\) in the \(n^{\text{th}}\) slot and \(0\) elsewhere. Then \(e_1, e_3, e_5, \ldots\) all lie in \(\operatorname{null} T^*\) and are linearly independent (indeed orthonormal). Hence \(\operatorname{null} T^*\) contains linearly independent lists of every length, so \(\dim \operatorname{null} T^* = \infty\).

Thus \(T\) is a bounded operator on the Hilbert space \(\ell^2\) with \(\dim \operatorname{null} T = 0 < \infty\) but \(\dim \operatorname{null} T^* = \infty\), giving the desired counterexample.

Remark. The conceptual reason is 10.13(a): \(\operatorname{null} T^* = (\operatorname{range} T)^{\perp}\). Here \(\operatorname{range} T\) is the closed subspace of elements supported on the even coordinates, whose orthogonal complement (the elements supported on the odd coordinates) is infinite-dimensional. Any injective bounded operator whose range has infinite codimension provides such a counterexample.

Problem (10A.6)

Suppose \(T\) is a bounded linear map from a Hilbert space \(V\) to a Hilbert space \(W\). Prove that \(\|T^*T\| = \|T\|^2\).

[This formula for \(\|T^*T\|\) leads to the important subject of \(C^*\)-algebras.]

Solution

Here \(T \in \mathcal{B}(V, W)\) and \(T^* \in \mathcal{B}(W, V)\) by 10.11, so \(T^*T \in \mathcal{B}(V)\) is a bounded operator on \(V\). We prove the two inequalities separately.

Step 1: \(\|T^*T\| \le \|T\|^2\).

By 10.20 (the norm of a composition) and by 10.11 (which gives \(\|T^*\| = \|T\|\)),

\begin{equation*} \|T^*T\| \le \|T^*\| \, \|T\| = \|T\| \, \|T\| = \|T\|^2 . \end{equation*}

Step 2: \(\|T\|^2 \le \|T^*T\|\).

Suppose \(f \in V\) with \(\|f\| \le 1\). Using the defining property 10.2 of the adjoint (with \(g = Tf\)), then the Cauchy–Schwarz inequality (8.11), then the definition of the operator norm, we get

\begin{equation*} \begin{aligned} \|Tf\|^2 &= \langle Tf, Tf \rangle = \langle f, T^*(Tf) \rangle = \langle f, (T^*T) f \rangle \\ &= |\langle f, (T^*T) f \rangle| \le \|f\| \, \|(T^*T) f\| \le \|f\| \, \|T^*T\| \, \|f\| = \|T^*T\| \, \|f\|^2 \le \|T^*T\| . \end{aligned} \end{equation*}

In the second line, replacing \(\langle f, (T^*T)f \rangle\) by its absolute value changes nothing because that number equals \(\|Tf\|^2\) and so is a nonnegative real number; the first inequality is then the Cauchy–Schwarz inequality (8.11), the second is \(\|(T^*T)f\| \le \|T^*T\| \, \|f\|\), and the last uses \(\|f\| \le 1\).

Hence \(\|Tf\| \le \|T^*T\|^{1/2}\) for every \(f \in V\) with \(\|f\| \le 1\). Taking the supremum over all such \(f\) and using

\begin{equation*} \|T\| = \sup \{ \|Tf\| : f \in V, \ \|f\| \le 1 \} \end{equation*}

gives \(\|T\| \le \|T^*T\|^{1/2}\), that is, \(\|T\|^2 \le \|T^*T\|\).

Combining Steps 1 and 2 yields

\begin{equation*} \|T^*T\| = \|T\|^2 , \end{equation*}

as desired.

Remark. Replacing \(T\) by \(T^*\) in the identity just proved and using \((T^*)^* = T\) and \(\|T^*\| = \|T\|\) (both from 10.11) gives the companion identity \(\|TT^*\| = \|T^*\|^2 = \|T\|^2\) as well. This identity \(\|T^*T\| = \|T\|^2\) is exactly the \(C^*\)-identity, the defining axiom relating the norm and the involution in a \(C^*\)-algebra.

Problem (10A.7)

Suppose \(V\) is a Hilbert space and \(\operatorname{Inv}(V)\) is the set of invertible bounded operators on \(V\). Think of \(\operatorname{Inv}(V)\) as a metric space with the metric it inherits as a subset of \(\mathcal{B}(V)\). Show that \(T \mapsto T^{-1}\) is a continuous function from \(\operatorname{Inv}(V)\) to \(\operatorname{Inv}(V)\).

Solution

Throughout, the metric on \(\mathcal{B}(V)\) is \(d(S, T) = \|S - T\|\), where \(\|\cdot\|\) is the operator norm.

Preliminaries. A Hilbert space is complete, so \(V\) is a Banach space. Hence if \(T \in \mathcal{B}(V)\) is invertible, then the inverse function \(T^{-1} : V \to V\) is linear and, by the Bounded Inverse Theorem (6.83), bounded. Thus \(T^{-1} \in \mathcal{B}(V)\), and \(T^{-1}\) is itself invertible (its inverse is \(T\)), so \(T^{-1} \in \operatorname{Inv}(V)\). Therefore \(T \mapsto T^{-1}\) genuinely maps \(\operatorname{Inv}(V)\) into \(\operatorname{Inv}(V)\).

If \(V = \{0\}\), then \(\mathcal{B}(V)\) contains only the zero operator and every function on \(\operatorname{Inv}(V)\) is continuous; so assume from now on that \(V \ne \{0\}\). Then for \(T \in \operatorname{Inv}(V)\) the operator \(T^{-1}\) is surjective onto \(V \ne \{0\}\), so \(T^{-1} \ne 0\) and hence \(\|T^{-1}\| > 0\).

The key algebraic identity. If \(S, T \in \operatorname{Inv}(V)\), then expanding the product and using \(TT^{-1} = I\) together with \(S^{-1}S = I\) gives

\begin{equation*} S^{-1}(T - S)T^{-1} = S^{-1}TT^{-1} - S^{-1}ST^{-1} = S^{-1} - T^{-1} . \end{equation*}

Hence, using 10.20 twice,

\begin{equation*} \|S^{-1} - T^{-1}\| \le \|S^{-1}\| \, \|T - S\| \, \|T^{-1}\| . \end{equation*}

So to prove continuity at \(T\) we need a bound on \(\|S^{-1}\|\) that is uniform for \(S\) near \(T\).

A uniform bound on \(\|S^{-1}\|\) near \(T\). Fix \(T \in \operatorname{Inv}(V)\), and suppose \(S \in \operatorname{Inv}(V)\) satisfies

\begin{equation*} \|T - S\| < \frac{1}{2\|T^{-1}\|} . \end{equation*}

Put \(A = I - T^{-1}S\). Then \(A = T^{-1}T - T^{-1}S = T^{-1}(T - S)\), so by 10.20,

\begin{equation*} \|A\| \le \|T^{-1}\| \, \|T - S\| < \tfrac{1}{2} . \end{equation*}

Since \(\|A\| < 1\), the operator \(I - A = T^{-1}S\) is invertible and \((T^{-1}S)^{-1} = \sum_{k=0}^{\infty} A^k\) by 10.22. Because the norm is continuous and satisfies the triangle inequality, and because \(\|A^k\| \le \|A\|^k\),

\begin{equation*} \|(T^{-1}S)^{-1}\| = \Big\| \sum_{k=0}^{\infty} A^k \Big\| \le \sum_{k=0}^{\infty} \|A\|^k = \frac{1}{1 - \|A\|} \le \frac{1}{1 - \frac12} = 2 . \end{equation*}

Now \(S^{-1} = S^{-1}(TT^{-1}) = (S^{-1}T)T^{-1} = (T^{-1}S)^{-1}T^{-1}\), so by 10.20 again,

\begin{equation*} \|S^{-1}\| \le \|(T^{-1}S)^{-1}\| \, \|T^{-1}\| \le 2\|T^{-1}\| . \end{equation*}

Continuity. Fix \(T \in \operatorname{Inv}(V)\) and let \(\varepsilon > 0\). Define

\begin{equation*} \delta = \min \left\{ \frac{1}{2\|T^{-1}\|}, \ \frac{\varepsilon}{2\|T^{-1}\|^2} \right\} > 0 . \end{equation*}

Suppose \(S \in \operatorname{Inv}(V)\) with \(\|S - T\| < \delta\). Then \(\|T - S\| < \frac{1}{2\|T^{-1}\|}\), so the previous paragraph gives \(\|S^{-1}\| \le 2\|T^{-1}\|\). Combining this with the key identity,

\begin{equation*} \|S^{-1} - T^{-1}\| \le \|S^{-1}\| \, \|T - S\| \, \|T^{-1}\| \le 2\|T^{-1}\|^2 \, \|T - S\| < 2\|T^{-1}\|^2 \delta \le \varepsilon . \end{equation*}

Thus for every \(\varepsilon > 0\) there exists \(\delta > 0\) such that \(S \in \operatorname{Inv}(V)\) and \(\|S - T\| < \delta\) imply \(\|S^{-1} - T^{-1}\| < \varepsilon\). Hence \(T \mapsto T^{-1}\) is continuous at \(T\). Because \(T \in \operatorname{Inv}(V)\) was arbitrary, the map \(T \mapsto T^{-1}\) is a continuous function from \(\operatorname{Inv}(V)\) to \(\operatorname{Inv}(V)\).

Remark. The estimate above in fact shows local Lipschitz continuity: on the ball of radius \(\frac{1}{2\|T^{-1}\|}\) about \(T\), inversion is Lipschitz with constant \(2\|T^{-1}\|^2\). Note also that by 10.25 every \(S \in \mathcal{B}(V)\) with \(\|T - S\| < \frac{1}{\|T^{-1}\|}\) is automatically invertible, so \(\operatorname{Inv}(V)\) is open in \(\mathcal{B}(V)\); hence inversion is a homeomorphism of the open set \(\operatorname{Inv}(V)\) onto itself, because it is a continuous bijection of \(\operatorname{Inv}(V)\) onto itself that is its own inverse function.

Problem (10A.8)

Suppose \(T\) is a bounded operator on a Hilbert space.

(a) Prove that \(T\) is left invertible if and only if \(T^*\) is right invertible.

(b) Prove that \(T\) is invertible if and only if \(T\) is both left and right invertible.

Solution

Let \(V\) denote the Hilbert space and let \(T \in \mathcal{B}(V)\).

(a) Suppose \(T\) is left invertible. Then there exists \(S \in \mathcal{B}(V)\) with \(ST = I\). Taking adjoints of both sides and using 10.12(d) and 10.12(c), we get

\begin{equation*} T^* S^* = (ST)^* = I^* = I . \end{equation*}

By 10.11 we have \(S^* \in \mathcal{B}(V)\), so the displayed equation exhibits \(S^*\) as a bounded operator with \(T^* S^* = I\). Thus \(T^*\) is right invertible.

Conversely, suppose \(T^*\) is right invertible, so there exists \(R \in \mathcal{B}(V)\) with \(T^* R = I\). Taking adjoints and using 10.12(d), 10.12(c), and \((T^*)^* = T\) (from 10.11), we get

\begin{equation*} R^* T = R^* (T^*)^* = (T^* R)^* = I^* = I . \end{equation*}

Because \(R^* \in \mathcal{B}(V)\) (by 10.11), this shows that \(T\) is left invertible.

(b) First suppose \(T\) is invertible. Then \(T^{-1} : V \to V\) is a linear map with \(T^{-1} T = T T^{-1} = I\), and \(T^{-1}\) is bounded by the Bounded Inverse Theorem (6.83), because \(V\) is a Banach space (every Hilbert space is complete). Hence \(T^{-1} \in \mathcal{B}(V)\) serves simultaneously as a left inverse and as a right inverse of \(T\), so \(T\) is both left invertible and right invertible.

Conversely, suppose \(T\) is both left invertible and right invertible. Choose \(S, R \in \mathcal{B}(V)\) with

\begin{equation*} ST = I \qquad \text{and} \qquad TR = I . \end{equation*}

Then, using associativity of composition,

\begin{equation*} S = SI = S(TR) = (ST)R = IR = R . \end{equation*}

Thus \(S = R\), and therefore \(ST = TS = I\). This means that \(T\) is invertible, with \(T^{-1} = S\). (Explicitly: \(ST = I\) forces \(T\) to be injective, because \(Tf = 0\) implies \(f = (ST)f = S(Tf) = 0\); and \(TS = I\) forces \(T\) to be surjective, because each \(g \in V\) equals \(T(Sg)\). Hence \(T\) is a one-to-one linear map of \(V\) onto \(V\), which is the definition 10.18 of invertibility.)

Problem (10A.9)

Suppose \(b_1, b_2, \ldots\) is a bounded sequence in \(\mathbf{F}\). Define a bounded linear map \(T : \ell^2 \to \ell^2\) by

\begin{equation*} T(a_1, a_2, \ldots) = (a_1 b_1, a_2 b_2, \ldots). \end{equation*}

(a) Find a formula for \(T^*\).

(b) Show that \(T\) is injective if and only if \(b_k \neq 0\) for every \(k \in \mathbf{Z}^+\).

(c) Show that \(T\) has dense range if and only if \(b_k \neq 0\) for every \(k \in \mathbf{Z}^+\).

(d) Show that \(T\) has closed range if and only if

\begin{equation*} \inf\{|b_k| : k \in \mathbf{Z}^+ \text{ and } b_k \neq 0\} > 0 . \end{equation*}

(e) Show that \(T\) is invertible if and only if \(\inf\{|b_k| : k \in \mathbf{Z}^+\} > 0\).

Solution

Throughout, \(e_k \in \ell^2\) denotes the sequence with \(1\) in slot \(k\) and \(0\) elsewhere. Put

\begin{equation*} N = \{k \in \mathbf{Z}^+ : b_k \neq 0\}, \qquad Z = \{k \in \mathbf{Z}^+ : b_k = 0\} . \end{equation*}

Note \(T e_k = b_k e_k\) for each \(k\). Also note that \(T\) is indeed bounded, with \(\|T\| \le \sup_k |b_k| < \infty\).

(a) For \(a = (a_1, a_2, \ldots)\) and \(c = (c_1, c_2, \ldots)\) in \(\ell^2\),

\begin{equation*} \langle Ta, c \rangle = \sum_{k=1}^{\infty} a_k b_k \overline{c_k} = \sum_{k=1}^{\infty} a_k \overline{\overline{b_k}\, c_k} = \langle a, (\overline{b_1} c_1, \overline{b_2} c_2, \ldots) \rangle \end{equation*}

(all series converge absolutely by the Cauchy-Schwarz inequality, since \((b_k c_k)\) lies in \(\ell^2\)). Hence

\begin{equation*} T^*(c_1, c_2, \ldots) = (\overline{b_1} c_1, \overline{b_2} c_2, \ldots) . \end{equation*}

In other words, \(T^*\) is the operator of the same form built from the bounded sequence \(\overline{b_1}, \overline{b_2}, \ldots\).

(b) If \(b_k = 0\) for some \(k\), then \(T e_k = 0\) while \(e_k \neq 0\), so \(T\) is not injective. Conversely, suppose \(b_k \neq 0\) for every \(k\) and \(Ta = 0\). Then \(a_k b_k = 0\) for every \(k\), which forces \(a_k = 0\) for every \(k\); thus \(a = 0\) and \(T\) is injective.

(c) By 10.14, \(T\) has dense range if and only if \(T^*\) is injective. By (a), \(T^*\) is the multiplier operator associated with the sequence \((\overline{b_k})\), so by (b) applied to \(T^*\), the map \(T^*\) is injective if and only if \(\overline{b_k} \neq 0\) for every \(k\), which happens if and only if \(b_k \neq 0\) for every \(k\).

(d) Write \(c = \inf\{|b_k| : k \in N\}\), with the convention \(\inf \emptyset = \infty\). Dispose first of the degenerate case \(N = \emptyset\): then \(b_k = 0\) for every \(k\), so \(T = 0\) and \(\operatorname{range} T = \{0\}\) is closed, while \(c = \infty > 0\); so the claimed equivalence holds. Assume from now on that \(N \neq \emptyset\), so \(c \in [0, \infty)\).

First suppose \(c > 0\). Set

\begin{equation*} W = \{a \in \ell^2 : a_k = 0 \text{ for every } k \in Z\} . \end{equation*}

Then \(W\) is a closed subspace of \(\ell^2\): if \(a^{(n)} \in W\) and \(a^{(n)} \to a\) in \(\ell^2\), then \(|a^{(n)}_k - a_k| \le \|a^{(n)} - a\|\) shows that each coordinate converges, so \(a_k = 0\) for \(k \in Z\). Clearly \(\operatorname{range} T \subseteq W\). Conversely, suppose \(a \in W\) and define \(d \in \mathbf{F}^{\mathbf{Z}^+}\) by \(d_k = a_k / b_k\) for \(k \in N\) and \(d_k = 0\) for \(k \in Z\). Then \(|d_k| \le |a_k| / c\) for every \(k\), so \(\|d\| \le \|a\| / c < \infty\), giving \(d \in \ell^2\); and \(Td = a\) because \(b_k d_k = a_k\) for \(k \in N\) and \(b_k d_k = 0 = a_k\) for \(k \in Z\). Hence \(\operatorname{range} T = W\), which is closed.

Now suppose \(c = 0\); we show \(\operatorname{range} T\) is not closed. For \(j \in \mathbf{Z}^+\) let \(N_j = \{k \in N : |b_k| < 1/j\}\). Because \(c = 0\), each \(N_j\) is nonempty. In fact each \(N_j\) is infinite: if some \(N_j\) were finite and nonempty, then \(\delta := \min\{|b_k| : k \in N_j\} > 0\); choosing \(j^{\prime}\) with \(1/j^{\prime} \le \delta\) and \(j^{\prime} \ge j\), any \(k \in N_{j^{\prime}}\) would satisfy \(|b_k| < 1/j^{\prime} \le \delta\) and \(k \in N_{j^{\prime}} \subseteq N_j\), hence \(|b_k| \ge \delta\), a contradiction; so \(N_{j^{\prime}} = \emptyset\), contradicting \(c = 0\). Thus we may choose \(k_1 < k_2 < \cdots\) with \(k_j \in N_j\), so that \(b_{k_j} \neq 0\) and \(|b_{k_j}| < 1/j\) for every \(j\).

Define \(a \in \ell^2\) by \(a_{k_j} = 1/j\) for each \(j\) and \(a_k = 0\) for all other \(k\); then \(\|a\|^2 = \sum_j 1/j^2 < \infty\). For each \(J \in \mathbf{Z}^+\) let \(a^{(J)}\) agree with \(a\) in slots \(k_1, \ldots, k_J\) and be \(0\) elsewhere. Since \(k_j \in N\), we have \(a^{(J)} = T d^{(J)}\) where \(d^{(J)}\) has \(j\)th relevant entry \(d^{(J)}_{k_j} = (1/j)/b_{k_j}\) for \(j \le J\) and \(0\) elsewhere; \(d^{(J)}\) is finitely supported, hence in \(\ell^2\). Thus \(a^{(J)} \in \operatorname{range} T\), and \(\|a - a^{(J)}\|^2 = \sum_{j > J} 1/j^2 \to 0\), so \(a\) lies in the closure of \(\operatorname{range} T\).

However \(a \notin \operatorname{range} T\). Indeed, if \(Td = a\) for some \(d \in \ell^2\), then \(d_{k_j} b_{k_j} = 1/j\), so

\begin{equation*} |d_{k_j}| = \frac{1}{j\,|b_{k_j}|} > \frac{1}{j \cdot (1/j)} = 1 \end{equation*}

for every \(j\), and hence \(\|d\|^2 \ge \sum_{j=1}^{\infty} |d_{k_j}|^2 = \infty\), contradicting \(d \in \ell^2\). Therefore \(\operatorname{range} T\) is not closed.

(e) Suppose \(\inf\{|b_k| : k \in \mathbf{Z}^+\} = c > 0\). Then \(b_k \neq 0\) for all \(k\), so \(Z = \emptyset\). Define \(S : \ell^2 \to \ell^2\) by \(S(a_1, a_2, \ldots) = (a_1/b_1, a_2/b_2, \ldots)\); this is well defined and bounded with \(\|S\| \le 1/c\), since \(|a_k/b_k| \le |a_k|/c\). Clearly \(ST = TS = I\), so \(T\) is invertible.

Conversely, suppose \(T\) is invertible. Then \(T^{-1}\) is a bounded operator (by 6.83, or by definition together with 10.18). For each \(k\) we have \(T e_k = b_k e_k\); in particular \(b_k \neq 0\), since otherwise \(T\) would not be injective. Applying \(T^{-1}\) gives \(T^{-1} e_k = e_k / b_k\), hence

\begin{equation*} \frac{1}{|b_k|} = \|T^{-1} e_k\| \le \|T^{-1}\| \,\|e_k\| = \|T^{-1}\| . \end{equation*}

Thus \(|b_k| \ge 1/\|T^{-1}\|\) for every \(k\) (note \(\|T^{-1}\| > 0\)), so \(\inf_k |b_k| \ge 1/\|T^{-1}\| > 0\).

Problem (10A.10)

Suppose \(h \in L^\infty(\mathbf{R})\) and \(M_h : L^2(\mathbf{R}) \to L^2(\mathbf{R})\) is the bounded operator defined by \(M_h f = fh\).

(a) Show that \(M_h\) is injective if and only if \(|\{x \in \mathbf{R} : h(x) = 0\}| = 0\).

(b) Find a necessary and sufficient condition (in terms of \(h\)) for \(M_h\) to have dense range.

(c) Find a necessary and sufficient condition (in terms of \(h\)) for \(M_h\) to have closed range.

(d) Find a necessary and sufficient condition (in terms of \(h\)) for \(M_h\) to be invertible.

Solution

Fix a representative of \(h\) and write \(Z = \{x \in \mathbf{R} : h(x) = 0\}\); changing the representative changes \(Z\) only by a set of measure \(0\), so all the conditions below are unambiguous. Here \(|A|\) denotes Lebesgue measure of \(A\). Note first that \(M_h\) is bounded, with \(\|M_h f\|_2 \le \|h\|_\infty \|f\|_2\).

We record the adjoint: for \(f, g \in L^2(\mathbf{R})\),

\begin{equation*} \langle M_h f, g \rangle = \int f h \overline{g} \, d\lambda = \int f \overline{\overline{h} g} \, d\lambda = \langle f, M_{\overline{h}} g \rangle , \end{equation*}

so \((M_h)^* = M_{\overline{h}}\).

(a) Suppose \(|Z| > 0\). Because \(Z = \bigcup_{n=1}^{\infty} (Z \cap [-n,n])\), the formula for the measure of an increasing union (2.59) gives some \(n\) with \(|Z \cap [-n,n]| > 0\); put \(E = Z \cap [-n,n]\), so \(0 < |E| < \infty\). Then \(\chi_E \in L^2(\mathbf{R})\) and \(\chi_E \neq 0\) in \(L^2\), while \(M_h \chi_E = h \chi_E = 0\) (because \(h\) vanishes on \(E\)). Hence \(M_h\) is not injective.

Conversely suppose \(|Z| = 0\) and \(M_h f = 0\). Then \(f h = 0\) almost everywhere, so for almost every \(x \notin Z\) we have \(f(x) h(x) = 0\) with \(h(x) \neq 0\), giving \(f(x) = 0\). Since \(|Z| = 0\), we conclude \(f = 0\) almost everywhere, i.e. \(f = 0\) in \(L^2(\mathbf{R})\). Thus \(M_h\) is injective.

(b) The condition is the same one as in (a): \(M_h\) has dense range if and only if \(|\{x \in \mathbf{R} : h(x) = 0\}| = 0\).

Indeed, by 10.14, \(M_h\) has dense range if and only if \((M_h)^* = M_{\overline{h}}\) is injective. By part (a) applied to \(\overline{h}\), this happens if and only if \(|\{x : \overline{h}(x) = 0\}| = 0\), and \(\{\overline{h} = 0\} = \{h = 0\} = Z\).

(c) The condition is: there exists \(\delta > 0\) such that

\begin{equation*} \bigl|\{x \in \mathbf{R} : 0 < |h(x)| < \delta\}\bigr| = 0 , \end{equation*}

i.e. \(|h|\) is essentially bounded away from \(0\) on the set where it does not vanish.

Suppose first that such a \(\delta > 0\) exists. Let

\begin{equation*} W = \{f \in L^2(\mathbf{R}) : f = 0 \text{ almost everywhere on } Z\} . \end{equation*}

Then \(W = \operatorname{null} M_{\chi_Z}\) is the null space of a bounded operator, hence a closed subspace of \(L^2(\mathbf{R})\). Clearly \(\operatorname{range} M_h \subseteq W\), since \(fh\) vanishes on \(Z\). Conversely let \(f \in W\) and define \(g = f/h\) on \(\mathbf{R} \setminus Z\) and \(g = 0\) on \(Z\). By hypothesis \(|h| \ge \delta\) almost everywhere on \(\mathbf{R} \setminus Z\), so \(|g| \le |f|/\delta\) almost everywhere; hence \(g \in L^2(\mathbf{R})\). Moreover \(gh = f\) almost everywhere: on \(\mathbf{R} \setminus Z\) this is the definition of \(g\), and on \(Z\) both sides vanish almost everywhere. Thus \(f = M_h g \in \operatorname{range} M_h\). Hence \(\operatorname{range} M_h = W\) is closed.

Now suppose no such \(\delta\) exists, so \(|\{0 < |h| < \delta\}| > 0\) for every \(\delta > 0\). Put

\begin{equation*} B_n = \Bigl\{ x \in \mathbf{R} : \tfrac{1}{n+1} \le |h(x)| < \tfrac{1}{n} \Bigr\} \qquad (n \in \mathbf{Z}^+), \end{equation*}

so the \(B_n\) are pairwise disjoint and \(\{0 < |h| < 1/n\} = \bigcup_{m \ge n} B_m\). If only finitely many \(B_n\) had positive measure, say \(|B_m| = 0\) for all \(m \ge M\), then countable additivity would give \(|\{0 < |h| < 1/M\}| = 0\), contrary to assumption. Hence there are \(n_1 < n_2 < \cdots\) with \(|B_{n_j}| > 0\) for every \(j\); note \(n_j \ge j\). Intersecting with a large interval as in (a), choose \(E_j \subseteq B_{n_j}\) with \(0 < |E_j| < \infty\). The sets \(E_1, E_2, \ldots\) are pairwise disjoint.

Define

\begin{equation*} f = \sum_{j=1}^{\infty} \frac{1}{j \, |E_j|^{1/2}} \, \chi_{E_j} . \end{equation*}

Since the \(E_j\) are disjoint, \(\|f\|_2^2 = \sum_{j} \frac{1}{j^2 |E_j|} \cdot |E_j| = \sum_j 1/j^2 < \infty\), so \(f \in L^2(\mathbf{R})\).

Let \(f^{(J)} = \sum_{j \le J} \frac{1}{j |E_j|^{1/2}} \chi_{E_j}\). On \(\bigcup_{j \le J} E_j\) we have \(|h| \ge \min_{j \le J} \frac{1}{n_j + 1} > 0\), so \(g^{(J)} := f^{(J)}/h\) (defined to be \(0\) off \(\bigcup_{j \le J} E_j\)) lies in \(L^2(\mathbf{R})\) and \(M_h g^{(J)} = f^{(J)}\). Since \(\|f - f^{(J)}\|_2^2 = \sum_{j > J} 1/j^2 \to 0\), the function \(f\) lies in the closure of \(\operatorname{range} M_h\).

But \(f \notin \operatorname{range} M_h\). Suppose \(M_h g = f\) for some \(g \in L^2(\mathbf{R})\). On \(E_j\) we have \(h \neq 0\), so \(g = f/h\) almost everywhere on \(E_j\), and there \(|h| < 1/n_j \le 1/j\), whence

\begin{equation*} |g| = \frac{|f|}{|h|} > \frac{1}{j\,|E_j|^{1/2}} \cdot j = \frac{1}{|E_j|^{1/2}} \end{equation*}

almost everywhere on \(E_j\). Therefore \(\int_{E_j} |g|^2 \, d\lambda \ge |E_j| / |E_j| = 1\) for every \(j\), and by disjointness \(\|g\|_2^2 \ge \sum_{j=1}^{\infty} 1 = \infty\), contradicting \(g \in L^2(\mathbf{R})\).

Hence \(\operatorname{range} M_h\) is not closed, completing the proof of the characterization.

(d) The condition is: there exists \(\delta > 0\) such that \(|\{x \in \mathbf{R} : |h(x)| < \delta\}| = 0\); equivalently, \(1/h \in L^\infty(\mathbf{R})\) (that is, \(|h|\) is essentially bounded away from \(0\) on all of \(\mathbf{R}\)).

Suppose such a \(\delta\) exists. Then \(|h| \ge \delta\) almost everywhere, so \(1/h\) (defined arbitrarily on the null set where \(h = 0\)) belongs to \(L^\infty(\mathbf{R})\) with \(\|1/h\|_\infty \le 1/\delta\). Hence \(M_{1/h}\) is a bounded operator on \(L^2(\mathbf{R})\), and \(M_{1/h} M_h = M_h M_{1/h} = M_{(1/h) h} = I\), because \((1/h)h = 1\) almost everywhere. Thus \(M_h\) is invertible.

Conversely suppose \(M_h\) is invertible. Because \(L^2(\mathbf{R})\) is a Banach space, the Bounded Inverse Theorem (6.83) guarantees that \((M_h)^{-1}\) is a bounded operator; let \(\alpha = \|(M_h)^{-1}\|\), and note \(\alpha > 0\) because \((M_h)^{-1} \neq 0\). For every \(f \in L^2(\mathbf{R})\),

\begin{equation*} \|f\|_2 = \|(M_h)^{-1} M_h f\|_2 \le \alpha \, \|M_h f\|_2 . \end{equation*}

Let \(\delta = 1/(2\alpha)\) and suppose, seeking a contradiction, that \(A := \{x : |h(x)| < \delta\}\) has positive measure. As in (a), choose \(E \subseteq A\) with \(0 < |E| < \infty\), and take \(f = \chi_E\). Then \(\|f\|_2 = |E|^{1/2}\), while

\begin{equation*} \|M_h f\|_2^2 = \int_E |h|^2 \, d\lambda \le \delta^2 |E| , \end{equation*}

so \(\|M_h f\|_2 \le \delta |E|^{1/2}\). The displayed inequality gives \(|E|^{1/2} \le \alpha \delta |E|^{1/2} = \tfrac12 |E|^{1/2}\), which is impossible because \(|E|^{1/2} > 0\). Hence \(|A| = 0\), i.e. \(|\{|h| < \delta\}| = 0\) with \(\delta = 1/(2\|(M_h)^{-1}\|)\).

Problem (10A.11)

(a) Prove or give a counterexample: If \(T\) is a bounded operator on a Hilbert space such that \(T\) and \(T^*\) are both injective, then \(T\) is invertible.

(b) Prove or give a counterexample: If \(T\) is a bounded operator on a Hilbert space such that \(T\) and \(T^*\) are both surjective, then \(T\) is invertible.

Solution

(a) This is false. Define \(T : \ell^2 \to \ell^2\) by

\begin{equation*} T(a_1, a_2, a_3, \ldots) = \Bigl(a_1, \frac{a_2}{2}, \frac{a_3}{3}, \ldots\Bigr) , \end{equation*}

which is the operator of Example 10.28; it is bounded because \(|a_k/k| \le |a_k|\), so \(\|Ta\| \le \|a\|\).

By the computation in Exercise 9(a) (with \(b_k = 1/k\), a real sequence), we have \(T^* = T\). So it suffices to show that \(T\) is injective but not invertible.

\(T\) is injective: if \(Ta = 0\) then \(a_k/k = 0\) for every \(k\), so \(a = 0\). Hence \(T\) and \(T^* = T\) are both injective.

\(T\) is not surjective: let \(a = (1, \tfrac12, \tfrac13, \ldots)\), so \(\|a\|^2 = \sum_k 1/k^2 < \infty\) and \(a \in \ell^2\). If \(Td = a\) for some \(d \in \ell^2\), then \(d_k / k = 1/k\), so \(d_k = 1\) for every \(k\), and \(d \notin \ell^2\). Hence \(a \notin \operatorname{range} T\), so \(T\) is not surjective and therefore not invertible.

(Equivalently: by Exercise 9(e), \(T\) is invertible if and only if \(\inf_k |b_k| > 0\), and here \(\inf_k 1/k = 0\).)

(b) This is true. Suppose \(T \in \mathcal{B}(V)\) for a Hilbert space \(V\), and suppose \(T\) and \(T^*\) are both surjective.

Because \(T^*\) is surjective, \(\operatorname{range} T^* = V\), and hence the closure of \(\operatorname{range} T^*\) also equals \(V\). By 10.13(b),

\begin{equation*} V = \overline{\operatorname{range} T^*} = (\operatorname{null} T)^\perp . \end{equation*}

Taking orthogonal complements and using 8.41 together with the fact that \(\operatorname{null} T\) is a closed subspace, we get

\begin{equation*} \operatorname{null} T = \bigl( (\operatorname{null} T)^\perp \bigr)^\perp = V^\perp = \{0\} . \end{equation*}

Thus \(T\) is injective. Since \(T\) is also surjective by hypothesis, \(T\) is a one-to-one linear map of \(V\) onto \(V\), hence invertible by 10.18.

(Alternatively: \(T\) surjective implies \(T\) is right invertible by 10.31, and \(T^*\) surjective implies (again by 10.31, applied to \(T^*\)) that \(T^*\) is right invertible, which by Exercise 8(a) means \(T\) is left invertible; now Exercise 8(b) gives that \(T\) is invertible.)

Problem (10A.12)

Define \(T : \ell^2 \to \ell^2\) by \(T(a_1, a_2, a_3, \ldots) = (a_2, a_3, a_4, \ldots)\). Suppose \(\alpha \in \mathbf{F}\).

(a) Prove that \(T - \alpha I\) is injective if and only if \(|\alpha| \ge 1\).

(b) Prove that \(T - \alpha I\) is invertible if and only if \(|\alpha| > 1\).

(c) Prove that \(T - \alpha I\) is surjective if and only if \(|\alpha| \neq 1\).

(d) Prove that \(T - \alpha I\) is left invertible if and only if \(|\alpha| > 1\).

Solution

Let \(S : \ell^2 \to \ell^2\) be the right shift, \(S(a_1, a_2, a_3, \ldots) = (0, a_1, a_2, \ldots)\). (These are the two operators of Example 10.27, where the right shift is called \(T\) and the left shift is called \(S\); here the names are interchanged so that \(T\) matches the operator of this exercise.) Three preliminary facts:

Fact 1. \(T^* = S\) and \(S^* = T\). Indeed, for \(a, c \in \ell^2\),

\begin{equation*} \langle Sa, c \rangle = \sum_{k=2}^{\infty} a_{k-1} \overline{c_k} = \sum_{j=1}^{\infty} a_j \overline{c_{j+1}} = \langle a, Tc \rangle , \end{equation*}

so \(S^* = T\); by 10.11, \(T^* = (S^*)^* = S\).

Fact 2. \(\|Ta\| \le \|a\|\) and \(\|Sa\| = \|a\|\) for all \(a \in \ell^2\), since \(\|Ta\|^2 = \sum_{k \ge 2} |a_k|^2 \le \|a\|^2\) and \(\|Sa\|^2 = \sum_{k \ge 1} |a_k|^2 = \|a\|^2\). In particular \(\|T\| \le 1\).

Fact 3. If \(|\alpha| = 1\), then \(T - \alpha I\) is not bounded below; that is, there is no \(c > 0\) with \(\|(T - \alpha I)a\| \ge c \|a\|\) for all \(a \in \ell^2\).

To see this, for \(n \in \mathbf{Z}^+\) let \(a^{(n)} = (1, \alpha, \alpha^2, \ldots, \alpha^{n-1}, 0, 0, \ldots)\), i.e. \(a^{(n)}_k = \alpha^{k-1}\) for \(k \le n\) and \(a^{(n)}_k = 0\) for \(k > n\). Since \(|\alpha| = 1\), we have \(\|a^{(n)}\| = \sqrt{n}\). The \(k\)th coordinate of \((T - \alpha I) a^{(n)}\) equals \(a^{(n)}_{k+1} - \alpha a^{(n)}_k\), which is \(\alpha^k - \alpha \cdot \alpha^{k-1} = 0\) for \(k \le n - 1\), equals \(0 - \alpha \cdot \alpha^{n-1} = -\alpha^n\) for \(k = n\), and equals \(0\) for \(k > n\). Hence \(\|(T - \alpha I)a^{(n)}\| = |\alpha|^n = 1\), so

\begin{equation*} \frac{\|(T - \alpha I)a^{(n)}\|}{\|a^{(n)}\|} = \frac{1}{\sqrt{n}} \longrightarrow 0 , \end{equation*}

proving Fact 3.

(a) The equation \((T - \alpha I)a = 0\) says precisely that \(a_{k+1} = \alpha a_k\) for every \(k \in \mathbf{Z}^+\), which by induction means \(a_k = \alpha^{k-1} a_1\) for every \(k\).

If \(|\alpha| < 1\), put \(a = (1, \alpha, \alpha^2, \ldots)\). Then \(\|a\|^2 = \sum_{k \ge 0} |\alpha|^{2k} = 1/(1 - |\alpha|^2) < \infty\), so \(a \in \ell^2\); and \(a \neq 0\) and \((T - \alpha I)a = 0\). Hence \(T - \alpha I\) is not injective.

If \(|\alpha| \ge 1\), suppose \(a \in \ell^2\) and \((T - \alpha I)a = 0\). If \(a_1 \neq 0\), then \(|a_k| = |\alpha|^{k-1} |a_1| \ge |a_1| > 0\) for all \(k\), so \(\sum_k |a_k|^2 = \infty\), contradicting \(a \in \ell^2\). Hence \(a_1 = 0\), and then \(a_k = \alpha^{k-1} a_1 = 0\) for every \(k\), so \(a = 0\). Thus \(T - \alpha I\) is injective.

(b) Suppose \(|\alpha| > 1\). Then \(\|T/\alpha\| \le \|T\|/|\alpha| \le 1/|\alpha| < 1\) by Fact 2, so 10.22 shows that \(I - T/\alpha\) is invertible (note that \(\ell^2\) is a Banach space). Hence

\begin{equation*} T - \alpha I = -\alpha\Bigl(I - \frac{T}{\alpha}\Bigr) \end{equation*}

is invertible, because a nonzero scalar multiple \(\beta A\) of an invertible operator \(A\) is invertible with inverse \(\beta^{-1} A^{-1}\).

Conversely, suppose \(T - \alpha I\) is invertible. Then \(T - \alpha I\) is injective, so \(|\alpha| \ge 1\) by (a). Suppose, seeking a contradiction, that \(|\alpha| = 1\). Because \(\ell^2\) is a Banach space, the Bounded Inverse Theorem (6.83) shows that \((T - \alpha I)^{-1}\) is a bounded operator; it is not the zero operator, so \(c := 1/\|(T - \alpha I)^{-1}\|\) satisfies \(c > 0\). For every \(a \in \ell^2\),

\begin{equation*} \|a\| = \bigl\|(T - \alpha I)^{-1} (T - \alpha I) a\bigr\| \le \frac{1}{c} \, \|(T - \alpha I) a\| , \end{equation*}

so \(\|(T - \alpha I)a\| \ge c\|a\|\) for all \(a \in \ell^2\), contradicting Fact 3. Hence \(|\alpha| > 1\).

(c) Suppose \(|\alpha| > 1\). Then \(T - \alpha I\) is invertible by (b), hence surjective.

Suppose \(|\alpha| < 1\). By Fact 1 and 10.12,

\begin{equation*} (T - \alpha I)^* = T^* - \overline{\alpha} I = S - \overline{\alpha} I . \end{equation*}

For every \(b \in \ell^2\), Fact 2 gives

\begin{equation*} \|(S - \overline{\alpha} I) b\| \ge \|Sb\| - |\alpha| \, \|b\| = (1 - |\alpha|) \|b\| , \end{equation*}

and \(1 - |\alpha| > 0\). Thus condition (b) of 10.29 holds for \(S - \overline{\alpha} I\) (with \(\beta = 1/(1 - |\alpha|)\)), so \(S - \overline{\alpha} I = (T - \alpha I)^*\) is left invertible. By Exercise 8(a) applied to the operator \((T - \alpha I)^*\), together with \(\bigl((T - \alpha I)^*\bigr)^* = T - \alpha I\) (from 10.11), the operator \(T - \alpha I\) is right invertible. By 10.31, \(T - \alpha I\) is surjective.

Suppose \(|\alpha| = 1\). Then \(T - \alpha I\) is injective by (a). If it were also surjective, it would be a one-to-one linear map of \(\ell^2\) onto \(\ell^2\), hence invertible by 10.18. But (b) says that \(T - \alpha I\) is invertible only when \(|\alpha| > 1\), whereas \(|\alpha| = 1\). This contradiction shows that \(T - \alpha I\) is not surjective.

Combining the three cases: \(T - \alpha I\) is surjective if and only if \(|\alpha| \neq 1\).

(d) Suppose \(|\alpha| > 1\). Then \(T - \alpha I\) is invertible by (b), and its inverse is bounded by 6.83, so \((T - \alpha I)^{-1}\) is a left inverse of \(T - \alpha I\) in \(\mathcal{B}(\ell^2)\). Hence \(T - \alpha I\) is left invertible.

Suppose \(|\alpha| < 1\). By (a), \(T - \alpha I\) is not injective, so condition (c) of 10.29 fails; by the equivalence of (a) and (c) in 10.29, the operator \(T - \alpha I\) is not left invertible.

Suppose \(|\alpha| = 1\). By Fact 3 there is no \(\beta \in (0, \infty)\) with \(\|a\| \le \beta \|(T - \alpha I)a\|\) for all \(a \in \ell^2\), so condition (b) of 10.29 fails; by the equivalence of (a) and (b) in 10.29, the operator \(T - \alpha I\) is not left invertible.

Thus \(T - \alpha I\) is left invertible if and only if \(|\alpha| > 1\).

Problem (10A.13)

Suppose \(V\) is a Hilbert space.

(a) Show that \(\{T \in \mathcal{B}(V) : T \text{ is left invertible}\}\) is an open subset of \(\mathcal{B}(V)\).

(b) Show that \(\{T \in \mathcal{B}(V) : T \text{ is right invertible}\}\) is an open subset of \(\mathcal{B}(V)\).

Solution

(a) Suppose \(T \in \mathcal{B}(V)\) is left invertible. By 10.29 (the equivalence of (a) and (b)), there exists \(\alpha \in (0, \infty)\) such that

\begin{equation*} \|f\| \le \alpha \|Tf\| \qquad \text{for all } f \in V . \end{equation*}

We claim that every \(R \in \mathcal{B}(V)\) with \(\|T - R\| < 1/\alpha\) is left invertible; this shows that the open ball of radius \(1/\alpha\) centered at \(T\) is contained in the set of left invertible operators, and hence that this set is open.

So suppose \(\|T - R\| < 1/\alpha\) and set \(\varepsilon = 1/\alpha - \|T - R\| > 0\). For every \(f \in V\),

\begin{equation*} \|Rf\| \ge \|Tf\| - \|(T - R)f\| \ge \frac{\|f\|}{\alpha} - \|T - R\|\, \|f\| = \varepsilon \|f\| , \end{equation*}

where the second inequality uses the displayed bound \(\|f\| \le \alpha\|Tf\|\) and the definition of the operator norm. Thus \(\|f\| \le (1/\varepsilon) \|Rf\|\) for all \(f \in V\), so condition (b) of 10.29 holds for \(R\) (with the constant \(1/\varepsilon\)). By 10.29, \(R\) is left invertible, proving the claim.

(b) By 10.12(a) and 10.12(b) (with \(\alpha = -1\)), the map \(A \mapsto A^*\) satisfies \((T - R)^* = T^* - R^*\), and by 10.11 we have \(\|(T-R)^*\| = \|T - R\|\). Hence

\begin{equation*} \|T^* - R^*\| = \|T - R\| \qquad \text{for all } T, R \in \mathcal{B}(V) . \end{equation*}

Now suppose \(T \in \mathcal{B}(V)\) is right invertible. By Exercise 8(a) applied to \(T^*\), using \((T^*)^* = T\) from 10.11, the operator \(T^*\) is left invertible. By part (a) there exists \(r > 0\) such that every \(A \in \mathcal{B}(V)\) with \(\|T^* - A\| < r\) is left invertible.

Suppose now \(R \in \mathcal{B}(V)\) and \(\|T - R\| < r\). Then \(\|T^* - R^*\| = \|T - R\| < r\), so \(R^*\) is left invertible. Applying Exercise 8(a) to \(R^*\) and using \((R^*)^* = R\), we conclude that \(R\) is right invertible.

Thus the open ball of radius \(r\) centered at \(T\) consists of right invertible operators, so the set of right invertible operators in \(\mathcal{B}(V)\) is open.

Problem (10A.14)

Suppose \(T\) is a bounded operator on a Hilbert space \(V\).

(a) Prove that \(T\) is invertible if and only if \(T\) has a unique left inverse. In other words, prove that \(T\) is invertible if and only if there exists a unique \(S \in \mathcal{B}(V)\) such that \(ST = I\).

(b) Prove that \(T\) is invertible if and only if \(T\) has a unique right inverse. In other words, prove that \(T\) is invertible if and only if there exists a unique \(S \in \mathcal{B}(V)\) such that \(TS = I\).

Solution

(a) First suppose \(T\) is invertible. Then \(T^{-1} \in \mathcal{B}(V)\) (by the Bounded Inverse Theorem 6.83, since \(V\) is a Banach space) and \(T^{-1} T = I\), so \(T\) has at least one left inverse. If \(S \in \mathcal{B}(V)\) satisfies \(ST = I\), then

\begin{equation*} S = S I = S(T T^{-1}) = (ST) T^{-1} = I T^{-1} = T^{-1} . \end{equation*}

Hence \(T^{-1}\) is the only left inverse of \(T\), so \(T\) has a unique left inverse.

Conversely, suppose there is exactly one \(S \in \mathcal{B}(V)\) with \(ST = I\). Define

\begin{equation*} S^{\prime} = S + (I - TS) . \end{equation*}

Then \(S^{\prime} \in \mathcal{B}(V)\), because \(\mathcal{B}(V)\) is closed under addition and under composition (10.20 shows \(TS\) is bounded). Moreover

\begin{equation*} S^{\prime} T = ST + (I - TS)T = I + T - T(ST) = I + T - TI = I , \end{equation*}

so \(S^{\prime}\) is also a left inverse of \(T\). By the assumed uniqueness, \(S^{\prime} = S\), which forces \(I - TS = 0\), i.e.

\begin{equation*} TS = I . \end{equation*}

Thus \(T\) is both left invertible (via \(ST = I\)) and right invertible (via \(TS = I\)), so \(T\) is invertible by Exercise 8(b).

(b) First suppose \(T\) is invertible. As above, \(T^{-1} \in \mathcal{B}(V)\) and \(T T^{-1} = I\), so \(T\) has at least one right inverse. If \(S \in \mathcal{B}(V)\) satisfies \(TS = I\), then

\begin{equation*} S = I S = (T^{-1} T) S = T^{-1} (TS) = T^{-1} I = T^{-1} , \end{equation*}

so \(T^{-1}\) is the only right inverse of \(T\).

Conversely, suppose there is exactly one \(S \in \mathcal{B}(V)\) with \(TS = I\). Define

\begin{equation*} S^{\prime} = S + (I - ST) \in \mathcal{B}(V) . \end{equation*}

Then

\begin{equation*} T S^{\prime} = TS + T(I - ST) = I + T - (TS)T = I + T - IT = I , \end{equation*}

so \(S^{\prime}\) is also a right inverse of \(T\). By uniqueness, \(S^{\prime} = S\), which forces \(I - ST = 0\), i.e. \(ST = I\). Thus \(T\) is both right invertible and left invertible, so \(T\) is invertible by Exercise 8(b).

Exercises 10B

Problem (10B.1)

Verify all the assertions in Example 10.33.

Example 10.33 (eigenvalues and spectrum) asserts the following.

  • Suppose \(b_1, b_2, \ldots\) is a bounded sequence in \(\mathbf{F}\). Define a bounded linear map \(T : \ell^2 \to \ell^2\) by

\begin{equation*} T(a_1, a_2, \ldots) = (a_1 b_1, a_2 b_2, \ldots). \end{equation*}

Then the set of eigenvalues of \(T\) equals \(\{b_k : k \in \mathbf{Z}^+\}\) and the spectrum of \(T\) equals the closure of \(\{b_k : k \in \mathbf{Z}^+\}\).

  • Suppose \(h \in L^\infty(\mathbf{R})\). Define a bounded linear map \(M_h : L^2(\mathbf{R}) \to L^2(\mathbf{R})\) by \(M_h f = f h\). Then \(\alpha \in \mathbf{F}\) is an eigenvalue of \(M_h\) if and only if \(|\{t \in \mathbf{R} : h(t) = \alpha\}| > 0\). Also, \(\alpha \in \operatorname{sp}(M_h)\) if and only if \(|\{t \in \mathbf{R} : |h(t) - \alpha| < \varepsilon\}| > 0\) for all \(\varepsilon > 0\).

  • Define the right shift \(T : \ell^2 \to \ell^2\) and the left shift \(S : \ell^2 \to \ell^2\) by

\begin{equation*} T(a_1, a_2, a_3, \ldots) = (0, a_1, a_2, a_3, \ldots) \quad \text{and} \quad S(a_1, a_2, a_3, \ldots) = (a_2, a_3, a_4, \ldots). \end{equation*}

Then \(T\) has no eigenvalues, and \(\operatorname{sp}(T) = \{\alpha \in \mathbf{F} : |\alpha| \le 1\}\). Also, the set of eigenvalues of \(S\) is the open set \(\{\alpha \in \mathbf{F} : |\alpha| < 1\}\), and the spectrum of \(S\) is the closed set \(\{\alpha \in \mathbf{F} : |\alpha| \le 1\}\).

Solution

Throughout, \(e_k\) denotes the element of \(\ell^2\) whose \(k^{\text{th}}\) coordinate is \(1\) and whose other coordinates are \(0\). We use repeatedly the elementary observation that an invertible bounded operator \(R\) on a Banach space is bounded below: if \(Rf = g\) then \(\|f\| = \|R^{-1}g\| \le \|R^{-1}\| \, \|Rf\|\), so

\begin{equation*} \|f\| \le \|R^{-1}\| \, \|Rf\| \quad \text{for all } f. \end{equation*}

Hence if there is a sequence of unit vectors \(f_1, f_2, \ldots\) with \(\lim_{n \to \infty} \|R f_n\| = 0\), then \(R\) is not invertible.

First bullet point. Let \(M = \sup_k |b_k| < \infty\). Then \(\|Ta\|^2 = \sum_k |b_k|^2 |a_k|^2 \le M^2 \|a\|^2\), so \(T\) is a bounded operator on \(\ell^2\) with \(\|T\| \le M\).

Fix \(\alpha \in \mathbf{F}\). Then \(T - \alpha I\) is again a multiplication operator of the same type, namely

\begin{equation*} (T - \alpha I)(a_1, a_2, \ldots) = \bigl(a_1(b_1 - \alpha), a_2(b_2 - \alpha), \ldots \bigr). \end{equation*}

Eigenvalues: \(T - \alpha I\) is not injective if and only if there is a nonzero \(a \in \ell^2\) with \(a_k(b_k - \alpha) = 0\) for every \(k\). If \(\alpha = b_j\) for some \(j\), then \(e_j\) is such a vector, so \(\alpha\) is an eigenvalue. Conversely, if \(\alpha \ne b_k\) for every \(k\), then \(a_k(b_k - \alpha) = 0\) forces \(a_k = 0\) for every \(k\), so \(T - \alpha I\) is injective. Thus the set of eigenvalues of \(T\) is exactly \(\{b_k : k \in \mathbf{Z}^+\}\).

Spectrum: we claim \(T - \alpha I\) is invertible if and only if \(\inf_k |b_k - \alpha| > 0\). Suppose \(c := \inf_k |b_k - \alpha| > 0\). Define \(S : \ell^2 \to \ell^2\) by

\begin{equation*} S(a_1, a_2, \ldots) = \Bigl( \frac{a_1}{b_1 - \alpha}, \frac{a_2}{b_2 - \alpha}, \ldots \Bigr). \end{equation*}

Then \(\|Sa\| \le \|a\|/c\), so \(S\) is a bounded operator, and clearly \(S(T - \alpha I) = (T - \alpha I)S = I\). Hence \(T - \alpha I\) is invertible. Conversely, suppose \(\inf_k |b_k - \alpha| = 0\). Choose \(k_1, k_2, \ldots\) with \(\lim_{n \to \infty} |b_{k_n} - \alpha| = 0\). Then \(\|e_{k_n}\| = 1\) and \(\|(T - \alpha I)e_{k_n}\| = |b_{k_n} - \alpha| \to 0\), so by the observation at the start, \(T - \alpha I\) is not invertible.

Finally, \(\inf_k |b_k - \alpha| = 0\) if and only if \(\alpha\) is a limit of a sequence from \(\{b_k : k \in \mathbf{Z}^+\}\), which is exactly the statement that \(\alpha\) lies in the closure of \(\{b_k : k \in \mathbf{Z}^+\}\). Thus \(\operatorname{sp}(T)\) equals the closure of \(\{b_k : k \in \mathbf{Z}^+\}\).

Second bullet point. Fix \(\alpha \in \mathbf{F}\) and set \(g = h - \alpha\), so \(g \in L^\infty(\mathbf{R})\) and \(M_h - \alpha I = M_g\).

Eigenvalues: suppose \(|\{t \in \mathbf{R} : g(t) = 0\}| > 0\). The set \(Z = \{t : g(t) = 0\}\) is measurable, and \(Z = \bigcup_{n=1}^\infty \bigl(Z \cap [-n, n]\bigr)\), so by continuity from below there exists \(n\) with \(0 < |Z \cap [-n,n]| < \infty\). Put \(E = Z \cap [-n,n]\). Then \(\chi_E \in L^2(\mathbf{R})\), \(\chi_E \ne 0\) in \(L^2(\mathbf{R})\), and \(M_g \chi_E = g \chi_E = 0\). Hence \(M_g\) is not injective, so \(\alpha\) is an eigenvalue of \(M_h\). Conversely, suppose \(|Z| = 0\). If \(f \in L^2(\mathbf{R})\) and \(M_g f = 0\), then \(gf = 0\) almost everywhere; since \(g \ne 0\) almost everywhere, \(f = 0\) almost everywhere. Thus \(M_g\) is injective and \(\alpha\) is not an eigenvalue. Therefore \(\alpha\) is an eigenvalue of \(M_h\) if and only if \(|\{t \in \mathbf{R} : h(t) = \alpha\}| > 0\).

Spectrum: suppose there exists \(\varepsilon > 0\) with \(|\{t : |g(t)| < \varepsilon\}| = 0\); in other words \(|g| \ge \varepsilon\) almost everywhere. Then \(1/g\) (defined arbitrarily, say as \(0\), on the null set where \(|g| < \varepsilon\)) belongs to \(L^\infty(\mathbf{R})\) with \(\|1/g\|_\infty \le 1/\varepsilon\), and \(M_{1/g}\) is a bounded operator on \(L^2(\mathbf{R})\) satisfying \(M_{1/g} M_g = M_g M_{1/g} = I\) (the products agree with the identity almost everywhere). Hence \(M_g\) is invertible and \(\alpha \notin \operatorname{sp}(M_h)\).

Conversely, suppose \(|\{t : |g(t)| < \varepsilon\}| > 0\) for every \(\varepsilon > 0\). Fix \(\varepsilon > 0\). Exactly as above, the set \(\{t : |g(t)| < \varepsilon\}\) contains a measurable subset \(F\) with \(0 < |F| < \infty\). Put \(f = \chi_F / |F|^{1/2}\). Then \(\|f\|_2 = 1\) and

\begin{equation*} \|M_g f\|_2^2 = \frac{1}{|F|} \int_F |g|^2 \, d\lambda \le \varepsilon^2 . \end{equation*}

Letting \(\varepsilon\) run through \(1, \tfrac12, \tfrac13, \ldots\) produces unit vectors \(f_1, f_2, \ldots\) with \(\|M_g f_n\|_2 \to 0\), so \(M_g\) is not invertible. Hence \(\alpha \in \operatorname{sp}(M_h)\). This proves that \(\alpha \in \operatorname{sp}(M_h)\) if and only if \(|\{t : |h(t) - \alpha| < \varepsilon\}| > 0\) for all \(\varepsilon > 0\).

Third bullet point. The right shift \(T\) satisfies \(\|Ta\| = \|a\|\) for all \(a \in \ell^2\), so \(\|T\| = 1\). The left shift \(S\) satisfies \(\|Sa\| \le \|a\|\) and \(\|Se_2\| = \|e_1\| = 1\), so \(\|S\| = 1\).

\(T\) has no eigenvalues: suppose \(Ta = \alpha a\) with \(a \in \ell^2\). Comparing first coordinates gives \(0 = \alpha a_1\), and comparing the \((k+1)^{\text{st}}\) coordinates for \(k \ge 1\) gives \(a_k = \alpha a_{k+1}\). If \(\alpha = 0\), the second family of equations gives \(a_k = 0\) for all \(k\). If \(\alpha \ne 0\), the first equation gives \(a_1 = 0\), and then \(a_k = \alpha a_{k+1}\) inductively gives \(a_{k+1} = a_k / \alpha = 0\) for all \(k\). Either way \(a = 0\), so \(T\) has no eigenvector and hence no eigenvalue.

\(\operatorname{sp}(T) = \{\alpha : |\alpha| \le 1\}\): by 10.34(a) and \(\|T\| = 1\) we have \(\operatorname{sp}(T) \subseteq \{\alpha \in \mathbf{F} : |\alpha| \le 1\}\). For the reverse inclusion, let \(\alpha \in \mathbf{F}\) with \(|\alpha| < 1\) and set

\begin{equation*} v = (1, \overline{\alpha}, \overline{\alpha}^2, \ldots), \end{equation*}

which lies in \(\ell^2\) because \(|\alpha| < 1\); note \(v \ne 0\). For any \(a \in \ell^2\) we have \((T - \alpha I)a = (-\alpha a_1, \ a_1 - \alpha a_2, \ a_2 - \alpha a_3, \ldots)\), and therefore

\begin{equation*} \begin{aligned} \langle (T - \alpha I)a, v \rangle &= -\alpha a_1 + \sum_{k=1}^\infty (a_k - \alpha a_{k+1}) \alpha^k \\ &= -\alpha a_1 + \sum_{k=1}^\infty a_k \alpha^k - \sum_{k=1}^\infty a_{k+1}\alpha^{k+1} \\ &= -\alpha a_1 + a_1 \alpha = 0, \end{aligned} \end{equation*}

where all series converge absolutely by the Cauchy–Schwarz inequality (8.11). Thus the range of \(T - \alpha I\) is contained in \(\{v\}^\perp \ne \ell^2\), so \(T - \alpha I\) is not surjective and hence not invertible. Therefore \(\{\alpha \in \mathbf{F} : |\alpha| < 1\} \subseteq \operatorname{sp}(T)\). Because \(\operatorname{sp}(T)\) is closed (10.36), we get \(\{\alpha \in \mathbf{F} : |\alpha| \le 1\} \subseteq \operatorname{sp}(T)\), and equality follows.

Eigenvalues of \(S\): if \(|\alpha| < 1\), then \(w = (1, \alpha, \alpha^2, \ldots) \in \ell^2\), \(w \ne 0\), and \(Sw = (\alpha, \alpha^2, \alpha^3, \ldots) = \alpha w\); hence \(\alpha\) is an eigenvalue of \(S\). Conversely, suppose \(Sa = \alpha a\) with \(a \in \ell^2\) and \(a \ne 0\). Comparing \(k^{\text{th}}\) coordinates gives \(a_{k+1} = \alpha a_k\) for all \(k \ge 1\), so \(a_k = \alpha^{k-1} a_1\). If \(a_1 = 0\) then \(a = 0\), a contradiction; hence \(a_1 \ne 0\). Then \(\sum_k |a_k|^2 = |a_1|^2 \sum_{k \ge 0} |\alpha|^{2k} < \infty\) forces \(|\alpha| < 1\). Thus the set of eigenvalues of \(S\) is exactly \(\{\alpha \in \mathbf{F} : |\alpha| < 1\}\).

Spectrum of \(S\): every eigenvalue lies in the spectrum, so \(\{\alpha : |\alpha| < 1\} \subseteq \operatorname{sp}(S)\); since \(\operatorname{sp}(S)\) is closed (10.36), \(\{\alpha : |\alpha| \le 1\} \subseteq \operatorname{sp}(S)\). The reverse inclusion follows from 10.34(a) and \(\|S\| = 1\). Hence \(\operatorname{sp}(S) = \{\alpha \in \mathbf{F} : |\alpha| \le 1\}\).

Problem (10B.2)

Suppose \(T\) is a bounded operator on a Hilbert space \(V\).

  • (a) Prove that \(\operatorname{sp}(S^{-1}TS) = \operatorname{sp}(T)\) for all bounded invertible operators \(S\) on \(V\).

  • (b) Prove that \(\operatorname{sp}(T^*) = \{\overline{\alpha} : \alpha \in \operatorname{sp}(T)\}\).

  • (c) Prove that if \(T\) is invertible, then \(\operatorname{sp}(T^{-1}) = \bigl\{ \tfrac{1}{\alpha} : \alpha \in \operatorname{sp}(T) \bigr\}\).

Solution

We use throughout the following elementary fact: if \(A, B, C\) are bounded operators on \(V\) with \(A\) and \(C\) invertible, then \(ABC\) is invertible if and only if \(B\) is invertible. Indeed, if \(B\) is invertible then \(C^{-1}B^{-1}A^{-1}\) is a two-sided inverse of \(ABC\); conversely, if \(R\) is a two-sided inverse of \(ABC\), then \(CRA\) is a two-sided inverse of \(B\), because \(B(CRA) = A^{-1}(ABC)R A = A^{-1}A = I\) and \((CRA)B = C R (ABC) C^{-1} = C C^{-1} = I\).

(a) Suppose \(S\) is a bounded invertible operator on \(V\) and \(\alpha \in \mathbf{F}\). Then

\begin{equation*} S^{-1}TS - \alpha I = S^{-1}TS - \alpha S^{-1}S = S^{-1}(T - \alpha I)S . \end{equation*}

By the fact above (with \(A = S^{-1}\), \(B = T - \alpha I\), \(C = S\)), the operator \(S^{-1}TS - \alpha I\) is invertible if and only if \(T - \alpha I\) is invertible. Hence \(\alpha \in \operatorname{sp}(S^{-1}TS)\) if and only if \(\alpha \in \operatorname{sp}(T)\), which gives \(\operatorname{sp}(S^{-1}TS) = \operatorname{sp}(T)\).

(b) First note that if \(R\) is a bounded operator on \(V\), then \(R\) is invertible if and only if \(R^*\) is invertible. Indeed, if \(R^{-1}\) exists then taking adjoints in \(R R^{-1} = R^{-1} R = I\) and using 10.12(c) and 10.12(d) gives \((R^{-1})^* R^* = R^* (R^{-1})^* = I\), so \(R^*\) is invertible with \((R^*)^{-1} = (R^{-1})^*\). Applying this implication to \(R^*\) in place of \(R\) and using \((R^*)^* = R\) (see 10.11) gives the converse.

Now fix \(\beta \in \mathbf{F}\). By 10.12(a), 10.12(b), and 10.12(c),

\begin{equation*} T^* - \beta I = T^* - \overline{\overline{\beta}} \, I^* = \bigl(T - \overline{\beta} I \bigr)^* . \end{equation*}

Hence \(T^* - \beta I\) is invertible if and only if \(T - \overline{\beta}I\) is invertible. Equivalently,

\begin{equation*} \beta \in \operatorname{sp}(T^*) \iff \overline{\beta} \in \operatorname{sp}(T) . \end{equation*}

Thus \(\operatorname{sp}(T^*) = \{\beta : \overline{\beta} \in \operatorname{sp}(T)\} = \{\overline{\alpha} : \alpha \in \operatorname{sp}(T)\}\), where the last equality holds because \(\beta \mapsto \overline{\beta}\) is an involution of \(\mathbf{F}\). (When \(\mathbf{F} = \mathbf{R}\) the complex conjugation is the identity and the statement reads \(\operatorname{sp}(T^*) = \operatorname{sp}(T)\).)

(c) Suppose \(T\) is invertible. Then \(0 \notin \operatorname{sp}(T)\), so \(1/\alpha\) makes sense for every \(\alpha \in \operatorname{sp}(T)\); likewise \(T^{-1}\) is invertible, so \(0 \notin \operatorname{sp}(T^{-1})\).

Fix \(\alpha \in \mathbf{F}\) with \(\alpha \ne 0\). Then

\begin{equation*} -\frac{1}{\alpha} T^{-1} (T - \alpha I) = -\frac{1}{\alpha}\bigl(I - \alpha T^{-1}\bigr) = T^{-1} - \frac{1}{\alpha} I . \end{equation*}

The operator \(-\tfrac{1}{\alpha}T^{-1}\) is invertible, so by the fact stated at the beginning (with \(C = I\)), the operator \(T^{-1} - \tfrac{1}{\alpha}I\) is invertible if and only if \(T - \alpha I\) is invertible. In other words, for every nonzero \(\alpha \in \mathbf{F}\),

\begin{equation*} \frac{1}{\alpha} \in \operatorname{sp}(T^{-1}) \iff \alpha \in \operatorname{sp}(T). \end{equation*}

This gives both inclusions. If \(\alpha \in \operatorname{sp}(T)\), then \(\alpha \ne 0\) and the displayed equivalence puts \(1/\alpha \in \operatorname{sp}(T^{-1})\); hence \(\{1/\alpha : \alpha \in \operatorname{sp}(T)\} \subseteq \operatorname{sp}(T^{-1})\). Conversely, if \(\beta \in \operatorname{sp}(T^{-1})\), then \(\beta \ne 0\) (as noted above), and applying the equivalence with \(\alpha = 1/\beta\) shows \(1/\beta \in \operatorname{sp}(T)\); thus \(\beta = 1/(1/\beta) \in \{1/\alpha : \alpha \in \operatorname{sp}(T)\}\). Therefore

\begin{equation*} \operatorname{sp}(T^{-1}) = \Bigl\{ \frac{1}{\alpha} : \alpha \in \operatorname{sp}(T) \Bigr\}. \end{equation*}

Problem (10B.3)

Suppose \(E\) is a bounded subset of \(\mathbf{F}\). Show that there exists a Hilbert space \(V\) and \(T \in \mathcal{B}(V)\) such that the set of eigenvalues of \(T\) equals \(E\).

Solution

If \(E = \varnothing\), take \(V = \{0\}\) and \(T = 0\). The only vector in \(V\) is \(0\), and eigenvectors are by definition nonzero, so \(T\) has no eigenvalues; equivalently, \(T - \alpha I = 0\) is injective on \(V = \{0\}\) for every \(\alpha \in \mathbf{F}\). Thus the set of eigenvalues of \(T\) equals \(\varnothing = E\).

Now suppose \(E \ne \varnothing\). Let \(V = \ell^2(E)\), the Hilbert space \(L^2(\mu)\) where \(\mu\) is counting measure on the set \(E\) (see the second bullet point of Example 8.51). Concretely, \(V\) consists of the functions \(f : E \to \mathbf{F}\) such that

\begin{equation*} \|f\|^2 = \sum_{\alpha \in E} |f(\alpha)|^2 < \infty , \end{equation*}

with inner product \(\langle f, g \rangle = \sum_{\alpha \in E} f(\alpha)\overline{g(\alpha)}\). Note that no measurability or countability hypothesis on \(E\) is needed: every subset of \(E\) is measurable with respect to counting measure, so \(E\) may even be a nonmeasurable subset of \(\mathbf{F}\).

Because \(E\) is bounded, there exists \(M \in [0, \infty)\) with \(|\alpha| \le M\) for all \(\alpha \in E\). Define \(T : V \to V\) by

\begin{equation*} (Tf)(\alpha) = \alpha f(\alpha) \qquad \text{for } \alpha \in E . \end{equation*}

If \(f \in V\), then

\begin{equation*} \|Tf\|^2 = \sum_{\alpha \in E} |\alpha|^2 |f(\alpha)|^2 \le M^2 \sum_{\alpha \in E} |f(\alpha)|^2 = M^2 \|f\|^2 , \end{equation*}

so \(Tf \in V\) and \(\|Tf\| \le M \|f\|\). Linearity of \(T\) is clear. Hence \(T \in \mathcal{B}(V)\).

Every element of \(E\) is an eigenvalue of \(T\). Indeed, fix \(\beta \in E\) and let \(e_\beta \in V\) be the function that equals \(1\) at \(\beta\) and \(0\) elsewhere; then \(\|e_\beta\| = 1\), so \(e_\beta \ne 0\), and for every \(\alpha \in E\) we have \((Te_\beta)(\alpha) = \alpha e_\beta(\alpha) = \beta e_\beta(\alpha)\), because both sides are \(0\) unless \(\alpha = \beta\). Thus \(Te_\beta = \beta e_\beta\), so \(\beta\) is an eigenvalue of \(T\) with eigenvector \(e_\beta\).

Conversely, every eigenvalue of \(T\) lies in \(E\). Suppose \(\beta \in \mathbf{F}\) is an eigenvalue of \(T\), so there exists \(f \in V\) with \(f \ne 0\) and \(Tf = \beta f\). Then \((\alpha - \beta)f(\alpha) = 0\) for every \(\alpha \in E\). Because \(f \ne 0\), there exists \(\alpha_0 \in E\) with \(f(\alpha_0) \ne 0\), and then \(\alpha_0 - \beta = 0\), so \(\beta = \alpha_0 \in E\).

Therefore the set of eigenvalues of \(T\) equals \(E\).

Problem (10B.4)

Suppose \(E\) is a nonempty closed bounded subset of \(\mathbf{F}\). Show that there exists \(T \in \mathcal{B}(\ell^2)\) such that \(\operatorname{sp}(T) = E\).

Solution

The key point is that \(E\), being a subset of the separable metric space \(\mathbf{F}\), contains a countable dense subset. We prove this first.

Let \(D\) be a countable dense subset of \(\mathbf{F}\); take \(D = \mathbf{Q}\) if \(\mathbf{F} = \mathbf{R}\) and \(D = \{p + qi : p, q \in \mathbf{Q}\}\) if \(\mathbf{F} = \mathbf{C}\). For each pair \((d, n) \in D \times \mathbf{Z}^+\) such that

\begin{equation*} \Bigl\{ x \in E : |x - d| < \tfrac{1}{n} \Bigr\} \ne \varnothing , \end{equation*}

choose (using the axiom of choice for this countable family of nonempty sets) a point \(x_{d,n}\) in that set. Let \(A\) be the set of all points chosen in this way. Then \(A \subseteq E\) and \(A\) is countable, being indexed by a subset of the countable set \(D \times \mathbf{Z}^+\). Also \(A \ne \varnothing\): pick any \(x \in E\) (possible since \(E \ne \varnothing\)) and any \(d \in D\) with \(|x - d| < 1\); then the pair \((d, 1)\) is one of the pairs used above.

The set \(A\) is dense in \(E\). To see this, let \(x \in E\) and \(\varepsilon > 0\). Choose \(n \in \mathbf{Z}^+\) with \(\tfrac{2}{n} < \varepsilon\), and choose \(d \in D\) with \(|x - d| < \tfrac1n\) (possible because \(D\) is dense in \(\mathbf{F}\)). The set \(\{y \in E : |y - d| < \tfrac1n\}\) contains \(x\) and hence is nonempty, so \(x_{d,n}\) was chosen and satisfies \(|x_{d,n} - d| < \tfrac1n\). Therefore

\begin{equation*} |x_{d,n} - x| \le |x_{d,n} - d| + |d - x| < \tfrac{2}{n} < \varepsilon . \end{equation*}

Thus every point of \(E\) is a limit of points of \(A\), so \(E \subseteq \overline{A}\). Since \(A \subseteq E\) and \(E\) is closed, \(\overline{A} \subseteq E\). Hence

\begin{equation*} \overline{A} = E . \end{equation*}

Because \(A\) is countable and nonempty, we may list its elements as a sequence \(b_1, b_2, \ldots\) (if \(A\) is finite, repeat elements so as to obtain an infinite list); then \(\{b_k : k \in \mathbf{Z}^+\} = A\). Because \(E\) is bounded, the sequence \(b_1, b_2, \ldots\) is a bounded sequence in \(\mathbf{F}\).

Now define \(T : \ell^2 \to \ell^2\) by

\begin{equation*} T(a_1, a_2, \ldots) = (a_1 b_1, a_2 b_2, \ldots). \end{equation*}

If \(M = \sup_k |b_k| < \infty\), then \(\|Ta\|^2 = \sum_k |b_k|^2 |a_k|^2 \le M^2 \|a\|^2\), so \(T\) is well defined and \(T \in \mathcal{B}(\ell^2)\) with \(\|T\| \le M\).

By the first bullet point of Example 10.33 (verified in Exercise 1 of this section),

\begin{equation*} \operatorname{sp}(T) = \overline{\{b_k : k \in \mathbf{Z}^+\}} = \overline{A} = E , \end{equation*}

as desired.

Problem (10B.5)

Give an example of a bounded operator \(T\) on a normed vector space such that for every \(\alpha \in \mathbf{F}\), the operator \(T - \alpha I\) is not invertible.

Solution

Let

\begin{equation*} V = \{ a \in \ell^2 : a_k = 0 \text{ for all but finitely many } k \in \mathbf{Z}^+ \}, \end{equation*}

the subspace of \(\ell^2\) consisting of the sequences with only finitely many nonzero coordinates, equipped with the norm it inherits from \(\ell^2\). Then \(V\) is a normed vector space. It is not a Banach space, and it could not be: 10.34(b) shows that if \(W\) is a Banach space and \(R\) is a bounded operator on \(W\), then \(R - \alpha I\) is invertible for every \(\alpha \in \mathbf{F}\) with \(|\alpha| > \|R\|\). So no bounded operator on a Banach space, real or complex, has the property requested here, and incompleteness of the norm is essential.

Define \(T : V \to V\) to be the right shift:

\begin{equation*} T(a_1, a_2, a_3, \ldots) = (0, a_1, a_2, a_3, \ldots). \end{equation*}

If \(a \in V\) has all coordinates past the \(n^{\text{th}}\) equal to \(0\), then \(Ta\) has all coordinates past the \((n+1)^{\text{st}}\) equal to \(0\); thus \(T\) indeed maps \(V\) into \(V\). Clearly \(T\) is linear, and \(\|Ta\| = \|a\|\) for every \(a \in V\), so \(T\) is a bounded operator on \(V\) with \(\|T\| = 1\).

Now fix \(\alpha \in \mathbf{F}\); we show that \(T - \alpha I\) is not surjective, which implies that \(T - \alpha I\) is not invertible (if \(S\) were an inverse, then \((T - \alpha I)Sg = g\) for every \(g \in V\), forcing surjectivity).

Let \(e_1 = (1, 0, 0, \ldots) \in V\) and suppose, for contradiction, that there exists \(a \in V\) with \((T - \alpha I)a = e_1\). Writing out coordinates,

\begin{equation*} (T - \alpha I)a = \bigl( -\alpha a_1,\ a_1 - \alpha a_2,\ a_2 - \alpha a_3,\ \ldots \bigr). \end{equation*}

Comparing first coordinates gives \(-\alpha a_1 = 1\); in particular \(\alpha \ne 0\) and \(a_1 = -1/\alpha\). Comparing the \((k+1)^{\text{st}}\) coordinates for \(k \ge 1\) gives \(a_k - \alpha a_{k+1} = 0\), that is,

\begin{equation*} a_{k+1} = \frac{a_k}{\alpha} . \end{equation*}

By induction, \(a_k = -1/\alpha^{k}\) for every \(k \in \mathbf{Z}^+\). Thus every coordinate of \(a\) is nonzero, contradicting \(a \in V\).

Hence \(e_1 \notin \operatorname{range}(T - \alpha I)\), so \(T - \alpha I\) is not invertible. Because \(\alpha \in \mathbf{F}\) was arbitrary, \(T - \alpha I\) fails to be invertible for every \(\alpha \in \mathbf{F}\); equivalently, \(\operatorname{sp}(T) = \mathbf{F}\).

Problem (10B.6)

Suppose \(T\) is a bounded operator on a complex nonzero Banach space \(V\).

  • (a) Prove that the function

\begin{equation*} \alpha \mapsto \varphi\bigl( (T - \alpha I)^{-1} f \bigr) \end{equation*}

is analytic on \(\mathbf{C} \setminus \operatorname{sp}(T)\) for every \(f \in V\) and every \(\varphi \in V^{\prime}\).

  • (b) Prove that \(\operatorname{sp}(T) \ne \varnothing\).
Solution

Note first that 10.22, 10.25, 10.34, and 10.36 are all stated for bounded operators on Banach spaces, so they are available here. In particular \(\operatorname{sp}(T)\) is a closed subset of \(\mathbf{C}\) by 10.36, so \(\mathbf{C} \setminus \operatorname{sp}(T)\) is open and it makes sense to ask whether a function on it is analytic.

(a) Fix \(f \in V\) and \(\varphi \in V^{\prime}\), and let \(\beta \in \mathbf{C} \setminus \operatorname{sp}(T)\). Then \(T - \beta I\) is invertible; write

\begin{equation*} R = (T - \beta I)^{-1} \in \mathcal{B}(V). \end{equation*}

Note \(R \ne 0\), so \(\|R\| > 0\). Suppose \(\alpha \in \mathbf{C}\) with \(|\alpha - \beta| < 1/\|R\|\). Then \(\|(\alpha - \beta)R\| < 1\), so 10.22 shows that \(I - (\alpha - \beta)R\) is invertible with

\begin{equation*} \bigl( I - (\alpha - \beta) R \bigr)^{-1} = \sum_{k=0}^\infty (\alpha - \beta)^k R^k , \end{equation*}

the series converging in the norm of \(\mathcal{B}(V)\). Furthermore,

\begin{equation*} (T - \beta I)\bigl( I - (\alpha - \beta)R \bigr) = (T - \beta I) - (\alpha - \beta)(T - \beta I)R = (T - \beta I) - (\alpha - \beta) I = T - \alpha I . \end{equation*}

The left side is the composition of two invertible operators, so \(T - \alpha I\) is invertible and, using \((AB)^{-1} = B^{-1}A^{-1}\),

\begin{equation*} (T - \alpha I)^{-1} = \bigl( I - (\alpha - \beta) R \bigr)^{-1} R = \sum_{k=0}^\infty (\alpha - \beta)^k R^{k+1}, \end{equation*}

again with convergence in the norm of \(\mathcal{B}(V)\) (the series converges absolutely because \(\sum_k |\alpha - \beta|^k \|R\|^{k+1} < \infty\)).

The map \(S \mapsto \varphi(Sf)\) is a bounded linear functional on \(\mathcal{B}(V)\), because \(|\varphi(Sf)| \le \|\varphi\| \, \|f\| \, \|S\|\) for every \(S \in \mathcal{B}(V)\). Applying this continuous functional to the norm-convergent series above allows us to interchange it with the sum, giving

\begin{equation*} \varphi\bigl( (T - \alpha I)^{-1} f \bigr) = \sum_{k=0}^\infty \varphi\bigl( R^{k+1} f \bigr) (\alpha - \beta)^k \end{equation*}

for all \(\alpha \in \mathbf{C}\) with \(|\alpha - \beta| < 1/\|R\|\). The coefficients satisfy \(|\varphi(R^{k+1}f)| \le \|\varphi\| \, \|f\| \, \|R\|^{k+1}\), so this power series has radius of convergence at least \(1/\|R\|\). Thus the function \(\alpha \mapsto \varphi\bigl((T - \alpha I)^{-1}f\bigr)\) is given by a convergent power series in \(\alpha - \beta\) on a neighbourhood of \(\beta\), hence is analytic at \(\beta\). Because \(\beta\) was an arbitrary point of the open set \(\mathbf{C} \setminus \operatorname{sp}(T)\), this function is analytic on \(\mathbf{C} \setminus \operatorname{sp}(T)\).

(b) Suppose, for contradiction, that \(\operatorname{sp}(T) = \varnothing\). Then \(T - \alpha I\) is invertible for every \(\alpha \in \mathbf{C}\); in particular \(T\) itself is invertible (take \(\alpha = 0\)).

Because \(V \ne \{0\}\), we may choose \(f \in V\) with \(f \ne 0\). Fix \(\varphi \in V^{\prime}\) and define \(g : \mathbf{C} \to \mathbf{C}\) by

\begin{equation*} g(\alpha) = \varphi\bigl( (T - \alpha I)^{-1} f \bigr). \end{equation*}

By part (a) with \(\operatorname{sp}(T) = \varnothing\), the function \(g\) is analytic on all of \(\mathbf{C}\); that is, \(g\) is an entire function.

For every \(\alpha \in \mathbf{C}\),

\begin{equation*} |g(\alpha)| \le \|\varphi\| \, \bigl\| (T - \alpha I)^{-1} f \bigr\| \le \|\varphi\| \, \|f\| \, \bigl\| (T - \alpha I)^{-1} \bigr\| . \end{equation*}

By 10.34(c), \(\lim_{|\alpha| \to \infty} \|(T - \alpha I)^{-1}\| = 0\). Hence \(\lim_{|\alpha| \to \infty} g(\alpha) = 0\). In particular there exists \(r > 0\) such that \(|g(\alpha)| \le 1\) whenever \(|\alpha| \ge r\); and \(g\), being continuous, is bounded on the compact set \(\{\alpha : |\alpha| \le r\}\). Thus \(g\) is a bounded entire function, so Liouville’s theorem implies that \(g\) is constant. Since \(g(\alpha) \to 0\) as \(|\alpha| \to \infty\), that constant is \(0\). Hence \(g \equiv 0\); in particular

\begin{equation*} \varphi\bigl( T^{-1} f \bigr) = g(0) = 0 . \end{equation*}

The functional \(\varphi \in V^{\prime}\) was arbitrary, so \(\varphi(T^{-1}f) = 0\) for every \(\varphi \in V^{\prime}\). Let \(h = T^{-1}f\). If \(h \ne 0\), then 6.72 provides \(\varphi \in V^{\prime}\) with \(\|\varphi\| = 1\) and \(\varphi(h) = \|h\| > 0\), contradicting the previous sentence. Hence \(h = 0\), and therefore

\begin{equation*} f = T(T^{-1}f) = Th = 0, \end{equation*}

contradicting our choice of \(f \ne 0\).

This contradiction shows that the assumption \(\operatorname{sp}(T) = \varnothing\) is untenable. Thus \(\operatorname{sp}(T) \ne \varnothing\).

Problem (10B.7)

Prove that if \(T\) is an operator on a Hilbert space \(V\) such that \(\langle Tf, g \rangle = \langle f, Tg \rangle\) for all \(f, g \in V\), then \(T\) is a bounded operator.

Solution

Here \(T\) is a linear map from \(V\) to \(V\); no continuity is assumed, and the point of the exercise is that the symmetry hypothesis forces boundedness. (This result is often called the Hellinger–Toeplitz theorem. It is what makes the remark following 10.44 legitimate: for an everywhere-defined operator on a Hilbert space, the condition \(\langle Tf, g\rangle = \langle f, Tg\rangle\) already implies boundedness, hence self-adjointness.)

A Hilbert space is by definition a Banach space, so both the domain and the codomain of \(T\) are Banach spaces, and the Closed Graph Theorem (6.85) applies: \(T\) is a bounded linear map if and only if

\begin{equation*} \operatorname{graph}(T) = \{ (f, Tf) : f \in V \} \end{equation*}

is a closed subspace of \(V \times V\), where \(V \times V\) carries the norm \(\|(f, g)\| = \max\{\|f\|, \|g\|\}\) of 6.84.

First, \(\operatorname{graph}(T)\) is a subspace of \(V \times V\). Indeed, \((0, T0) = (0,0) \in \operatorname{graph}(T)\), and if \(f_1, f_2 \in V\) and \(\lambda \in \mathbf{F}\), then linearity of \(T\) gives

\begin{equation*} (f_1, Tf_1) + \lambda (f_2, Tf_2) = \bigl( f_1 + \lambda f_2, \ T(f_1 + \lambda f_2) \bigr) \in \operatorname{graph}(T). \end{equation*}

Now we show \(\operatorname{graph}(T)\) is closed. Suppose \((f_1, Tf_1), (f_2, Tf_2), \ldots\) is a sequence in \(\operatorname{graph}(T)\) converging to some \((f, g) \in V \times V\). By 6.84 this means

\begin{equation*} \lim_{k \to \infty} f_k = f \quad \text{and} \quad \lim_{k \to \infty} T f_k = g . \end{equation*}

Let \(h \in V\) be arbitrary. The Cauchy–Schwarz inequality (8.11) gives

\begin{equation*} \bigl| \langle T f_k, h \rangle - \langle g, h \rangle \bigr| = \bigl| \langle Tf_k - g, h \rangle \bigr| \le \| Tf_k - g \| \, \|h\| \longrightarrow 0 \end{equation*}

and

\begin{equation*} \bigl| \langle f_k, Th \rangle - \langle f, Th \rangle \bigr| = \bigl| \langle f_k - f, Th \rangle \bigr| \le \| f_k - f \| \, \|Th\| \longrightarrow 0 . \end{equation*}

Hence, using the hypothesis \(\langle Tf_k, h \rangle = \langle f_k, Th \rangle\) for each \(k\),

\begin{equation*} \langle g, h \rangle = \lim_{k \to \infty} \langle T f_k, h \rangle = \lim_{k \to \infty} \langle f_k, T h \rangle = \langle f, Th \rangle . \end{equation*}

Applying the hypothesis once more, this time with the pair \((f, h)\), gives \(\langle f, Th \rangle = \langle Tf, h \rangle\). Therefore

\begin{equation*} \langle g - Tf, h \rangle = 0 \qquad \text{for every } h \in V . \end{equation*}

Taking \(h = g - Tf\) yields \(\|g - Tf\|^2 = 0\), so \(g = Tf\). Thus \((f, g) = (f, Tf) \in \operatorname{graph}(T)\), proving that \(\operatorname{graph}(T)\) is closed.

Because \(\operatorname{graph}(T)\) is a closed subspace of \(V \times V\), the Closed Graph Theorem (6.85) implies that \(T\) is a bounded linear map, as desired.

Problem (10B.8)

Suppose \(P\) is a bounded operator on a Hilbert space \(V\) such that \(P^2 = P\). Prove that \(P\) is self-adjoint if and only if there exists a closed subspace \(U\) of \(V\) such that \(P = P_U\).

Solution

First suppose there is a closed subspace \(U\) of \(V\) with \(P = P_U\). The last bullet point of 10.45 shows that the orthogonal projection \(P_U\) onto a closed subspace is self-adjoint. Hence \(P\) is self-adjoint. (Note that this direction does not use the hypothesis \(P^2 = P\), which in any case holds automatically for \(P_U\).)

Now suppose \(P\) is self-adjoint. Set

\begin{equation*} U = \operatorname{range} P . \end{equation*}

We first show that \(U = \operatorname{null}(I - P)\). If \(f \in U\), then \(f = Pg\) for some \(g \in V\), and thus

\begin{equation*} P f = P^2 g = P g = f , \end{equation*}

which shows that \((I - P) f = 0\). Conversely, if \((I - P)f = 0\), then \(f = Pf \in \operatorname{range} P = U\). Thus indeed \(U = \operatorname{null}(I - P)\).

Because \(I - P\) is a bounded (hence continuous) linear map, its null space is a closed subspace of \(V\). Thus \(U\) is a closed subspace of \(V\), and \(P_U\) makes sense.

We now show \(P = P_U\). Let \(f \in V\). Certainly \(Pf \in \operatorname{range} P = U\). Suppose \(g \in U\). Because \(g \in U = \operatorname{null}(I - P)\) we have \(Pg = g\), and thus

\begin{equation*} \langle f - Pf , g \rangle = \langle f , g \rangle - \langle P f , g \rangle = \langle f , g \rangle - \langle f , P g \rangle = \langle f , g \rangle - \langle f , g \rangle = 0 , \end{equation*}

where the second equality holds because \(P\) is self-adjoint. Hence \(f - Pf\) is orthogonal to every element of \(U\), and \(Pf \in U\). Now 8.37(b) (applied with \(h = Pf\)) implies that

\begin{equation*} P f = P_U f . \end{equation*}

Because \(f \in V\) was arbitrary, \(P = P_U\), with \(U\) a closed subspace of \(V\), completing the proof.

Problem (10B.9)

Suppose \(V\) is a real Hilbert space and \(T \in \mathcal{B}(V)\). The complexification of \(T\) is the function \(T_{\mathbb{C}} : V_{\mathbb{C}} \to V_{\mathbb{C}}\) defined by

\begin{equation*} T_{\mathbb{C}}(f + ig) = Tf + iTg \end{equation*}

for \(f, g \in V\) (see Exercise 4 in Section 8B for the definition of \(V_{\mathbb{C}}\)).

(a) Show that \(T_{\mathbb{C}}\) is a bounded operator on the complex Hilbert space \(V_{\mathbb{C}}\) and \(\lVert T_{\mathbb{C}} \rVert = \lVert T \rVert\).

(b) Show that \(T_{\mathbb{C}}\) is invertible if and only if \(T\) is invertible.

(c) Show that \((T_{\mathbb{C}})^{*} = (T^{*})_{\mathbb{C}}\).

(d) Show that \(T\) is self-adjoint if and only if \(T_{\mathbb{C}}\) is self-adjoint.

(e) Use the previous parts of this exercise and 10.49 and 10.38 to show that if \(T\) is self-adjoint and \(V \neq \{0\}\), then \(\operatorname{sp}(T) \neq \emptyset\).

Solution

Throughout, recall from Exercise 4 in Section 8B that \(V_{\mathbb{C}} = V \times V\) with typical element written \(f + ig\), with

\begin{equation*} (\alpha + i\beta)(f + ig) = (\alpha f - \beta g) + i(\alpha g + \beta f) \end{equation*}

for \(\alpha, \beta \in \mathbb{R}\), and with inner product

\begin{equation*} \langle f_1 + i g_1 , f_2 + i g_2 \rangle = \langle f_1 , f_2 \rangle + \langle g_1 , g_2 \rangle + i \bigl( \langle g_1 , f_2 \rangle - \langle f_1 , g_2 \rangle \bigr) . \end{equation*}

Taking \(f_1 + ig_1 = f_2 + ig_2 = f + ig\) in the formula above and using the symmetry of the real inner product on \(V\) (so that \(\langle g , f \rangle - \langle f , g \rangle = 0\)), we obtain the identity that will be used repeatedly:

\begin{equation*} \lVert f + ig \rVert^2 = \lVert f \rVert^2 + \lVert g \rVert^2 . \end{equation*}

(a) First we check that \(T_{\mathbb{C}}\) is linear over \(\mathbb{C}\). Additivity is immediate from the additivity of \(T\):

\begin{equation*} T_{\mathbb{C}}\bigl( (f_1 + ig_1) + (f_2 + ig_2) \bigr) = T(f_1 + f_2) + i T(g_1 + g_2) = T_{\mathbb{C}}(f_1 + ig_1) + T_{\mathbb{C}}(f_2 + ig_2) . \end{equation*}

For homogeneity, let \(\alpha, \beta \in \mathbb{R}\) and \(f, g \in V\). Then

\begin{equation*} \begin{aligned} T_{\mathbb{C}}\bigl( (\alpha + i\beta)(f + ig) \bigr) &= T_{\mathbb{C}}\bigl( (\alpha f - \beta g) + i (\alpha g + \beta f) \bigr) \\ &= T(\alpha f - \beta g) + i T(\alpha g + \beta f) \\ &= (\alpha \, Tf - \beta \, Tg) + i (\alpha \, Tg + \beta \, Tf) \\ &= (\alpha + i \beta)(Tf + i Tg) \\ &= (\alpha + i\beta) \, T_{\mathbb{C}}(f + ig) , \end{aligned} \end{equation*}

where the third equality uses the \(\mathbb{R}\)-linearity of \(T\). Thus \(T_{\mathbb{C}}\) is a linear map on the complex vector space \(V_{\mathbb{C}}\).

For the norm, let \(f, g \in V\). Then

\begin{equation*} \lVert T_{\mathbb{C}}(f + ig) \rVert^2 = \lVert Tf \rVert^2 + \lVert Tg \rVert^2 \leq \lVert T \rVert^2 \bigl( \lVert f \rVert^2 + \lVert g \rVert^2 \bigr) = \lVert T \rVert^2 \, \lVert f + ig \rVert^2 . \end{equation*}

Hence \(T_{\mathbb{C}}\) is bounded with \(\lVert T_{\mathbb{C}} \rVert \leq \lVert T \rVert\). In the other direction, for \(f \in V\) we have \(\lVert f + i0 \rVert = \lVert f \rVert\) and

\begin{equation*} \lVert T_{\mathbb{C}}(f + i0) \rVert = \lVert Tf + i 0 \rVert = \lVert Tf \rVert . \end{equation*}

Taking the supremum over \(f \in V\) with \(\lVert f \rVert \leq 1\) gives \(\lVert T_{\mathbb{C}} \rVert \geq \lVert T \rVert\). Thus \(\lVert T_{\mathbb{C}} \rVert = \lVert T \rVert\).

(b) Suppose \(T\) is invertible, and let \(S = T^{-1} \in \mathcal{B}(V)\). By (a), \(S_{\mathbb{C}} \in \mathcal{B}(V_{\mathbb{C}})\). For all \(f, g \in V\),

\begin{equation*} S_{\mathbb{C}} T_{\mathbb{C}} (f + ig) = S_{\mathbb{C}}(Tf + iTg) = STf + i STg = f + ig , \end{equation*}

and symmetrically \(T_{\mathbb{C}} S_{\mathbb{C}} (f + ig) = f + ig\). Thus \(T_{\mathbb{C}}\) is invertible, with \((T_{\mathbb{C}})^{-1} = (T^{-1})_{\mathbb{C}}\).

Conversely, suppose \(T_{\mathbb{C}}\) is invertible, and let \(S = (T_{\mathbb{C}})^{-1} \in \mathcal{B}(V_{\mathbb{C}})\).

\(T\) is injective: if \(f \in V\) and \(Tf = 0\), then \(T_{\mathbb{C}}(f + i0) = 0 + i0\), so \(f + i0 = 0 + i0\) and hence \(f = 0\).

\(T\) is surjective: let \(f \in V\) and write \(S(f + i0) = u + iv\) with \(u, v \in V\). Applying \(T_{\mathbb{C}}\) gives \(Tu + iTv = f + i0\), so \(Tu = f\) and \(Tv = 0\). Thus \(f \in \operatorname{range} T\); moreover \(v = 0\) by injectivity, so \(S(f + i0) = u + i0\) where \(u = T^{-1}f\).

Hence \(T\) is a bijection, and its set-theoretic inverse \(T^{-1}\) is linear (the inverse of a linear bijection is linear). Finally \(T^{-1}\) is bounded: by the previous paragraph, for \(f \in V\),

\begin{equation*} \lVert T^{-1} f \rVert = \lVert u \rVert = \lVert u + i0 \rVert = \lVert S(f + i0) \rVert \leq \lVert S \rVert \, \lVert f + i0 \rVert = \lVert S \rVert \, \lVert f \rVert . \end{equation*}

Thus \(T\) is invertible.

(c) By (a), \((T^{*})_{\mathbb{C}}\) is a bounded operator on \(V_{\mathbb{C}}\). Let \(f_1, g_1, f_2, g_2 \in V\). Using the inner product formula on \(V_{\mathbb{C}}\),

\begin{equation*} \begin{aligned} \bigl\langle T_{\mathbb{C}}(f_1 + i g_1) , f_2 + i g_2 \bigr\rangle &= \langle Tf_1 + i Tg_1 , f_2 + i g_2 \rangle \\ &= \langle Tf_1 , f_2 \rangle + \langle Tg_1 , g_2 \rangle + i \bigl( \langle Tg_1 , f_2 \rangle - \langle Tf_1 , g_2 \rangle \bigr) . \end{aligned} \end{equation*}

On the other hand,

\begin{equation*} \begin{aligned} \bigl\langle f_1 + i g_1 , (T^{*})_{\mathbb{C}}(f_2 + i g_2) \bigr\rangle &= \langle f_1 + i g_1 , T^{*}f_2 + i T^{*} g_2 \rangle \\ &= \langle f_1 , T^{*} f_2 \rangle + \langle g_1 , T^{*} g_2 \rangle + i \bigl( \langle g_1 , T^{*} f_2 \rangle - \langle f_1 , T^{*} g_2 \rangle \bigr) \\ &= \langle Tf_1 , f_2 \rangle + \langle Tg_1 , g_2 \rangle + i \bigl( \langle Tg_1 , f_2 \rangle - \langle Tf_1 , g_2 \rangle \bigr) , \end{aligned} \end{equation*}

where the last equality uses \(\langle u , T^{*} w \rangle = \langle Tu , w \rangle\) for \(u , w \in V\) (the defining property of the adjoint on \(V\); the inner product on \(V\) is real and symmetric).

Thus \(\langle T_{\mathbb{C}} u , w \rangle = \langle u , (T^{*})_{\mathbb{C}} w \rangle\) for all \(u, w \in V_{\mathbb{C}}\). Because the adjoint of \(T_{\mathbb{C}}\) is the unique bounded operator on \(V_{\mathbb{C}}\) with this property, we conclude that \((T_{\mathbb{C}})^{*} = (T^{*})_{\mathbb{C}}\).

(d) The map \(S \mapsto S_{\mathbb{C}}\) is injective on \(\mathcal{B}(V)\): if \(S_{\mathbb{C}} = R_{\mathbb{C}}\), then evaluating at \(f + i0\) gives \(Sf = Rf\) for all \(f \in V\), so \(S = R\).

Now, using (c),

\begin{equation*} T \text{ is self-adjoint} \iff T = T^{*} \iff T_{\mathbb{C}} = (T^{*})_{\mathbb{C}} = (T_{\mathbb{C}})^{*} \iff T_{\mathbb{C}} \text{ is self-adjoint} , \end{equation*}

where the middle equivalence uses the injectivity just noted (for the direction from right to left) and the obvious implication (for the direction from left to right).

(e) Suppose \(T\) is self-adjoint and \(V \neq \{0\}\). Then \(V_{\mathbb{C}} \neq \{0\}\), because \(V_{\mathbb{C}}\) contains \(f + i0\) for each \(f \in V\) and \(\lVert f + i0 \rVert = \lVert f \rVert\). By Exercise 4 in Section 8B, \(V_{\mathbb{C}}\) is a complex Hilbert space, and by (a) we have \(T_{\mathbb{C}} \in \mathcal{B}(V_{\mathbb{C}})\). By (d), \(T_{\mathbb{C}}\) is self-adjoint.

Because \(V_{\mathbb{C}}\) is a nonzero complex Hilbert space, 10.38 implies that \(\operatorname{sp}(T_{\mathbb{C}}) \neq \emptyset\). Because \(T_{\mathbb{C}}\) is self-adjoint, 10.49 implies \(\operatorname{sp}(T_{\mathbb{C}}) \subseteq \mathbb{R}\). Choose \(\alpha \in \operatorname{sp}(T_{\mathbb{C}})\); thus \(\alpha \in \mathbb{R}\) and \(T_{\mathbb{C}} - \alpha I\) is not invertible.

Because \(\alpha\) is real, for all \(f, g \in V\) we have

\begin{equation*} (T - \alpha I)_{\mathbb{C}} (f + ig) = (Tf - \alpha f) + i (Tg - \alpha g) = T_{\mathbb{C}}(f + ig) - \alpha (f + ig) = (T_{\mathbb{C}} - \alpha I)(f + ig) , \end{equation*}

so \((T - \alpha I)_{\mathbb{C}} = T_{\mathbb{C}} - \alpha I\). Because \(T_{\mathbb{C}} - \alpha I\) is not invertible, (b) applied to the operator \(T - \alpha I \in \mathcal{B}(V)\) shows that \(T - \alpha I\) is not invertible. Since \(\alpha \in \mathbb{R} = \mathbb{F}\) (here \(\mathbb{F} = \mathbb{R}\) because \(V\) is a real Hilbert space), we conclude that \(\alpha \in \operatorname{sp}(T)\). Hence \(\operatorname{sp}(T) \neq \emptyset\).

Problem (10B.10)

Suppose \(T\) is a bounded operator on a Hilbert space \(V\) such that \(\langle Tf , f \rangle \geq 0\) for all \(f \in V\). Prove that \(\operatorname{sp}(T) \subseteq [0, \infty)\).

Solution

We show two things: that \(\operatorname{sp}(T) \subseteq \mathbb{R}\), and that no negative real number lies in \(\operatorname{sp}(T)\). Together these give \(\operatorname{sp}(T) \subseteq [0,\infty)\).

Step 1: we claim \(\operatorname{sp}(T) \subseteq \mathbb{R}\).

If \(\mathbb{F} = \mathbb{R}\), this is immediate from the definition 10.32, since \(\operatorname{sp}(T) \subseteq \mathbb{F} = \mathbb{R}\).

If \(\mathbb{F} = \mathbb{C}\), then the hypothesis says in particular that \(\langle Tf , f \rangle \in \mathbb{R}\) for every \(f \in V\); hence 10.48 implies that \(T\) is self-adjoint, and then 10.49 implies that \(\operatorname{sp}(T) \subseteq \mathbb{R}\).

Step 2: the adjoint satisfies the same hypothesis.

For \(f \in V\) we have

\begin{equation*} \langle T^{*} f , f \rangle = \overline{\langle f , T^{*} f \rangle} = \overline{\langle Tf , f \rangle} = \langle Tf , f \rangle \geq 0 , \end{equation*}

where the last equality holds because \(\langle Tf, f\rangle\) is a nonnegative real number. Thus \(T^{*}\) also satisfies the hypothesis of the exercise.

Step 3: we claim that if \(c > 0\) then \(T + cI\) is invertible.

Let \(c \in (0, \infty)\) and let \(f \in V\). Then

\begin{equation*} \langle (T + cI) f , f \rangle = \langle Tf , f \rangle + c \lVert f \rVert^2 \geq c \lVert f \rVert^2 \geq 0 . \end{equation*}

In particular \(\bigl| \langle (T + cI)f , f \rangle \bigr| \geq c \lVert f \rVert^2\). Now the Cauchy-Schwarz inequality (8.11) gives

\begin{equation*} \lVert (T + cI) f \rVert \, \lVert f \rVert \geq \bigl| \langle (T + cI) f , f \rangle \bigr| \geq c \lVert f \rVert^2 , \end{equation*}

and hence

\begin{equation*} \lVert (T + cI) f \rVert \geq c \lVert f \rVert \end{equation*}

for all \(f \in V\) (this is trivially true when \(f = 0\), and follows from the displayed inequality after dividing by \(\lVert f \rVert\) when \(f \neq 0\)). Equivalently, \(\lVert f \rVert \leq \frac{1}{c} \lVert (T+cI) f \rVert\) for all \(f \in V\). Thus the equivalence of (b) and (a) in 10.29 shows that \(T + cI\) is left invertible; say \(R (T + cI) = I\) for some \(R \in \mathcal{B}(V)\).

By Step 2, the operator \(T^{*}\) satisfies the same hypothesis as \(T\), so the argument just given (applied to \(T^{*}\) instead of \(T\)) shows that \(T^{*} + cI\) is left invertible: there exists \(R^{\prime} \in \mathcal{B}(V)\) with \(R^{\prime}(T^{*} + cI) = I\). Note that \(T^{*} + cI = (T + cI)^{*}\), since \(c\) is real. Taking adjoints in the equation \(R^{\prime}(T + cI)^{*} = I\) gives

\begin{equation*} (T + cI) (R^{\prime})^{*} = I , \end{equation*}

so \(T + cI\) is right invertible; set \(S = (R^{\prime})^{*}\).

Now \(T + cI\) is both left and right invertible, hence invertible: indeed

\begin{equation*} R = R I = R (T + cI) S = \bigl( R(T + cI) \bigr) S = I S = S , \end{equation*}

so \(R = S\) is a two-sided inverse of \(T + cI\) in \(\mathcal{B}(V)\).

Conclusion.

Let \(\alpha \in \operatorname{sp}(T)\). By Step 1, \(\alpha \in \mathbb{R}\). If \(\alpha < 0\), put \(c = -\alpha > 0\); then \(T - \alpha I = T + cI\) is invertible by Step 3, contradicting \(\alpha \in \operatorname{sp}(T)\). Hence \(\alpha \geq 0\). Therefore \(\operatorname{sp}(T) \subseteq [0, \infty)\).

Problem (10B.11)

Suppose \(P\) is a bounded operator on a Hilbert space \(V\) such that \(P^2 = P\). Prove that \(P\) is self-adjoint if and only if \(P\) is normal.

Solution

One direction is immediate: every self-adjoint operator is normal, because if \(P^{*} = P\) then \(P^{*}P = P^2 = PP^{*}\). (This holds for every bounded operator, with no need for \(P^2 = P\).)

For the other direction, suppose \(P\) is normal and \(P^2 = P\). We show that \(P = P_U\) for a closed subspace \(U\) of \(V\); then 10.45 (last bullet point) implies that \(P\) is self-adjoint. Equivalently, one may invoke Exercise 8 in this section at the last step.

As in Exercise 8, set \(U = \operatorname{range} P\). If \(f \in U\), say \(f = Pg\), then \(Pf = P^2 g = Pg = f\); conversely if \(Pf = f\) then \(f \in \operatorname{range} P\). Thus

\begin{equation*} U = \operatorname{range} P = \operatorname{null}(I - P) , \end{equation*}

which is closed because \(I - P\) is a bounded, hence continuous, linear map. So \(U\) is a closed subspace of \(V\).

Next we identify \(U^{\perp}\). For \(g \in V\) we have

\begin{equation*} g \in (\operatorname{range} P)^{\perp} \iff \langle Pf , g \rangle = 0 \text{ for all } f \in V \iff \langle f , P^{*} g \rangle = 0 \text{ for all } f \in V \iff P^{*} g = 0 , \end{equation*}

where the last equivalence holds by taking \(f = P^{*}g\). Thus

\begin{equation*} U^{\perp} = \operatorname{null} P^{*} . \end{equation*}

Now we use normality. By 10.53, since \(P\) is normal we have \(\lVert P h \rVert = \lVert P^{*} h \rVert\) for every \(h \in V\). Consequently \(Ph = 0\) if and only if \(P^{*}h = 0\), that is,

\begin{equation*} \operatorname{null} P = \operatorname{null} P^{*} = U^{\perp} . \end{equation*}

Finally, let \(f \in V\) and write

\begin{equation*} f = Pf + (f - Pf) . \end{equation*}

We have \(Pf \in \operatorname{range} P = U\). Also

\begin{equation*} P(f - Pf) = Pf - P^2 f = Pf - Pf = 0 , \end{equation*}

so \(f - Pf \in \operatorname{null} P = U^{\perp}\); in particular \(f - Pf\) is orthogonal to \(g\) for every \(g \in U\). Since \(Pf \in U\), 8.37(b) (with \(h = Pf\)) gives \(Pf = P_U f\).

Because \(f \in V\) was arbitrary, \(P = P_U\) with \(U\) a closed subspace of \(V\). By the last bullet point of 10.45, \(P_U\) is self-adjoint. Hence \(P\) is self-adjoint, completing the proof.

Problem (10B.12)

Prove that a normal operator on a separable Hilbert space has at most countably many eigenvalues.

Solution

Suppose \(V\) is a separable Hilbert space and \(T \in \mathcal{B}(V)\) is normal. Let

\begin{equation*} E = \{ \alpha \in \mathbb{F} : \alpha \text{ is an eigenvalue of } T \} . \end{equation*}

We must show that \(E\) is at most countable.

For each \(\alpha \in E\), choose an eigenvector of \(T\) corresponding to \(\alpha\) and divide it by its norm; this produces \(f_{\alpha} \in V\) with

\begin{equation*} \lVert f_{\alpha} \rVert = 1 \quad \text{and} \quad T f_{\alpha} = \alpha f_{\alpha} . \end{equation*}

By 10.57, eigenvectors of a normal operator corresponding to distinct eigenvalues are orthogonal. Thus if \(\alpha, \beta \in E\) with \(\alpha \neq \beta\), then \(\langle f_{\alpha} , f_{\beta} \rangle = 0\), and hence by the Pythagorean Theorem (8.9)

\begin{equation*} \lVert f_{\alpha} - f_{\beta} \rVert^2 = \lVert f_{\alpha} \rVert^2 + \lVert f_{\beta} \rVert^2 = 2 , \end{equation*}

so that

\begin{equation*} \lVert f_{\alpha} - f_{\beta} \rVert = \sqrt{2} . \end{equation*}

Because \(V\) is separable, by 8.64 there is a countable set \(D \subseteq V\) whose closure is \(V\). Hence for each \(\alpha \in E\) we may choose \(d_{\alpha} \in D\) with

\begin{equation*} \lVert f_{\alpha} - d_{\alpha} \rVert < \frac{\sqrt{2}}{2} . \end{equation*}

The function \(\alpha \mapsto d_{\alpha}\) from \(E\) to \(D\) is injective. Indeed, suppose \(\alpha, \beta \in E\) and \(d_{\alpha} = d_{\beta}\). Then

\begin{equation*} \lVert f_{\alpha} - f_{\beta} \rVert \leq \lVert f_{\alpha} - d_{\alpha} \rVert + \lVert d_{\beta} - f_{\beta} \rVert < \frac{\sqrt{2}}{2} + \frac{\sqrt{2}}{2} = \sqrt{2} . \end{equation*}

If \(\alpha \neq \beta\), this contradicts the equation \(\lVert f_{\alpha} - f_{\beta} \rVert = \sqrt{2}\) established above. Hence \(\alpha = \beta\).

Thus \(E\) injects into the countable set \(D\), which implies that \(E\) is at most countable. In other words, a normal operator on a separable Hilbert space has at most countably many eigenvalues.

Problem (10B.13)

Prove or give a counterexample: If \(T\) is a normal operator on a Hilbert space and \(T = A + iB\), where \(A\) and \(B\) are self-adjoint, then \(\lVert T \rVert = \sqrt{\lVert A \rVert^2 + \lVert B \rVert^2}\).

Solution

The statement is false. Here is a counterexample.

Let \(V = \mathbb{C}^2\) with the usual inner product, and define \(T \in \mathcal{B}(V)\) by

\begin{equation*} T(z_1 , z_2) = (z_1 , i z_2) . \end{equation*}

Thus the matrix of \(T\) with respect to the standard basis is the diagonal matrix with entries \(1\) and \(i\).

The adjoint of \(T\) is the operator whose matrix is the conjugate transpose (see 10.45), namely

\begin{equation*} T^{*}(z_1 , z_2) = (z_1 , -i z_2) . \end{equation*}

Then for all \((z_1, z_2) \in \mathbb{C}^2\),

\begin{equation*} T^{*} T (z_1 , z_2) = T^{*}(z_1 , i z_2) = (z_1 , z_2) \quad \text{and} \quad T T^{*}(z_1 , z_2) = T(z_1 , -i z_2) = (z_1 , z_2) . \end{equation*}

Hence \(T^{*}T = TT^{*} = I\), so \(T\) is normal (indeed unitary).

Now define \(A , B \in \mathcal{B}(V)\) by

\begin{equation*} A(z_1 , z_2) = (z_1 , 0) \quad \text{and} \quad B(z_1 , z_2) = (0 , z_2) . \end{equation*}

These are the operators produced by the formulas in 10.54, since

\begin{equation*} \frac{T + T^{*}}{2}(z_1 , z_2) = \frac{(z_1 , i z_2) + (z_1 , -i z_2)}{2} = (z_1 , 0) = A(z_1, z_2) \end{equation*}

and

\begin{equation*} \frac{T - T^{*}}{2i}(z_1 , z_2) = \frac{(0 , 2 i z_2)}{2i} = (0 , z_2) = B(z_1, z_2) . \end{equation*}

Both \(A\) and \(B\) are orthogonal projections onto closed subspaces of \(\mathbb{C}^2\) (onto the first and second coordinate axes respectively), so both are self-adjoint by the last bullet point of 10.45; this can also be seen directly from their matrices, each of which equals its own conjugate transpose. Furthermore

\begin{equation*} (A + iB)(z_1 , z_2) = (z_1 , 0) + i(0 , z_2) = (z_1 , i z_2) = T(z_1 , z_2) , \end{equation*}

so \(T = A + iB\) as required.

Now we compute the three norms. Because \(T^{*}T = I\), the operator \(T\) is an isometry (by the equivalence of (a) and (c) in 10.60), so \(\lVert T(z_1, z_2) \rVert = \lVert (z_1, z_2) \rVert\) for all \((z_1,z_2)\), and since \(V \neq \{0\}\) this gives

\begin{equation*} \lVert T \rVert = 1 . \end{equation*}

Also \(A(1,0) = (1,0)\) and \(\lVert A(z_1,z_2) \rVert = |z_1| \leq \lVert (z_1,z_2) \rVert\), so \(\lVert A \rVert = 1\); symmetrically \(B(0,1) = (0,1)\) and \(\lVert B(z_1,z_2) \rVert = |z_2| \leq \lVert (z_1,z_2) \rVert\), so \(\lVert B \rVert = 1\).

Therefore

\begin{equation*} \sqrt{\lVert A \rVert^2 + \lVert B \rVert^2} = \sqrt{1 + 1} = \sqrt{2} \neq 1 = \lVert T \rVert , \end{equation*}

which shows that the asserted equality fails.

A remark on why no counterexample can be one-dimensional: if \(V = \mathbb{C}\), then \(T\) is multiplication by some \(z = a + bi\) with \(a , b \in \mathbb{R}\), and then \(A\) and \(B\) are multiplication by \(a\) and by \(b\), so \(\lVert T \rVert = |z| = \sqrt{a^2 + b^2} = \sqrt{\lVert A \rVert^2 + \lVert B \rVert^2}\). The failure in higher dimensions occurs because \(A\) and \(B\) can attain their norms on different vectors, as they do in the example above. What does remain true in general is the inequality \(\max\{\lVert A \rVert , \lVert B \rVert\} \leq \lVert T \rVert\), which follows from \(\lVert A \rVert = \lVert \tfrac{1}{2}(T + T^{*}) \rVert \leq \tfrac{1}{2}(\lVert T \rVert + \lVert T^{*} \rVert) = \lVert T \rVert\) and the analogous estimate for \(B\), using \(\lVert T^{*} \rVert = \lVert T \rVert\) (by 10.11).

Problem (10B.14)

A number \(\alpha \in \mathbb{F}\) is called an approximate eigenvalue of a bounded operator \(T\) on a Hilbert space \(V\) if

\begin{equation*} \inf \bigl\{ \lVert (T - \alpha I) f \rVert : f \in V \text{ and } \lVert f \rVert = 1 \bigr\} = 0 . \end{equation*}

Suppose \(T\) is a normal operator on a Hilbert space and \(\alpha \in \mathbb{F}\). Prove that \(\alpha \in \operatorname{sp}(T)\) if and only if \(\alpha\) is an approximate eigenvalue of \(T\).

Solution

Let \(V\) denote the Hilbert space on which \(T\) acts. Dispose first of the degenerate case \(V = \{0\}\): then the only operator on \(V\) is \(0\), which is invertible, so \(\operatorname{sp}(T) = \emptyset\); and there are no vectors \(f \in V\) with \(\lVert f \rVert = 1\), so the infimum above is taken over the empty set and equals \(\infty \neq 0\), meaning \(\alpha\) is not an approximate eigenvalue. Thus both conditions fail for every \(\alpha\) and the equivalence holds. Assume from now on that \(V \neq \{0\}\).

First note that \(T - \alpha I\) is normal. Indeed, \((T - \alpha I)^{*} = T^{*} - \overline{\alpha} I\), and therefore

\begin{equation*} (T - \alpha I)^{*}(T - \alpha I) = T^{*}T - \alpha T^{*} - \overline{\alpha} T + |\alpha|^2 I \end{equation*}

and

\begin{equation*} (T - \alpha I)(T - \alpha I)^{*} = T T^{*} - \overline{\alpha} T - \alpha T^{*} + |\alpha|^2 I . \end{equation*}

The two right sides are equal because \(T^{*} T = T T^{*}\). Hence \(T - \alpha I\) is normal.

Suppose first that \(\alpha\) is not an approximate eigenvalue of \(T\). Then

\begin{equation*} c := \inf \bigl\{ \lVert (T - \alpha I) f \rVert : f \in V \text{ and } \lVert f \rVert = 1 \bigr\} > 0 . \end{equation*}

We claim that

\begin{equation*} \lVert f \rVert \leq \frac{1}{c} \lVert (T - \alpha I) f \rVert \qquad \text{for all } f \in V . \end{equation*}

This is trivial for \(f = 0\). If \(f \neq 0\), then \(f / \lVert f \rVert\) is a unit vector, so \(\bigl\lVert (T - \alpha I)(f / \lVert f \rVert) \bigr\rVert \geq c\); multiplying by \(\lVert f \rVert\) gives \(\lVert (T - \alpha I) f \rVert \geq c \lVert f \rVert\), which is the claim.

Thus condition (b) of 10.29 holds for the operator \(T - \alpha I\) (with the constant \(1/c\)), and hence by the equivalence of (a) and (b) in 10.29, the operator \(T - \alpha I\) is left invertible. Because \(T - \alpha I\) is normal, the equivalence of (b) and (a) in 10.55 implies that \(T - \alpha I\) is invertible. Hence \(\alpha \notin \operatorname{sp}(T)\).

Conversely, suppose \(\alpha \notin \operatorname{sp}(T)\), so that \(T - \alpha I\) is invertible; let \(S = (T - \alpha I)^{-1} \in \mathcal{B}(V)\). Because \(V \neq \{0\}\) and \(S\) is invertible, \(S \neq 0\), so \(\lVert S \rVert > 0\). If \(f \in V\) with \(\lVert f \rVert = 1\), then

\begin{equation*} 1 = \lVert f \rVert = \bigl\lVert S (T - \alpha I) f \bigr\rVert \leq \lVert S \rVert \, \bigl\lVert (T - \alpha I) f \bigr\rVert , \end{equation*}

and therefore

\begin{equation*} \bigl\lVert (T - \alpha I) f \bigr\rVert \geq \frac{1}{\lVert S \rVert} . \end{equation*}

Hence

\begin{equation*} \inf \bigl\{ \lVert (T - \alpha I) f \rVert : f \in V \text{ and } \lVert f \rVert = 1 \bigr\} \geq \frac{1}{\lVert S \rVert} > 0 , \end{equation*}

so \(\alpha\) is not an approximate eigenvalue of \(T\). (Note that this direction did not use the normality of \(T\); it holds for every bounded operator.)

Combining the two implications (each in contrapositive form) gives: \(\alpha \in \operatorname{sp}(T)\) if and only if \(\alpha\) is an approximate eigenvalue of \(T\).

It is worth noting where normality was essential. Without it, the first implication can fail: the right shift \(T\) on \(\ell^2\) satisfies \(\lVert T f \rVert = \lVert f \rVert\) for all \(f\), so \(0\) is not an approximate eigenvalue of \(T\), yet \(0 \in \operatorname{sp}(T)\) because \(T\) is not surjective (see 10.33 and 10.65). Normality is exactly what upgrades left invertibility to invertibility, via 10.55.

Problem (10B.15)

Suppose \(T\) is a normal operator on a Hilbert space.

(a) Prove that if \(\alpha\) is an eigenvalue of \(T\), then \(|\alpha|^2\) is an eigenvalue of \(T^*T\).

(b) Prove that if \(\alpha \in \operatorname{sp}(T)\), then \(|\alpha|^2 \in \operatorname{sp}(T^*T)\).

Solution

Throughout, \(V\) denotes the Hilbert space on which \(T\) acts.

(a) Suppose \(\alpha\) is an eigenvalue of \(T\), with eigenvector \(f \in V\), so \(f \neq 0\) and \(Tf = \alpha f\). Because \(T\) is normal, 10.56 tells us that \(f\) is also an eigenvector of \(T^*\), with eigenvalue \(\overline{\alpha}\); that is, \(T^* f = \overline{\alpha} f\). Hence

\begin{equation*} T^*Tf = T^*(\alpha f) = \alpha\, T^* f = \alpha \overline{\alpha} f = |\alpha|^2 f . \end{equation*}

Because \(f \neq 0\), this says exactly that \(|\alpha|^2\) is an eigenvalue of \(T^*T\), with eigenvector \(f\).

(b) Recall the terminology introduced just before Exercise 14 of this section: \(\beta \in \mathbf{F}\) is an approximate eigenvalue of a bounded operator \(S\) on \(V\) if

\begin{equation*} \inf\{\|(S - \beta I)f\| : f \in V \text{ and } \|f\| = 1\} = 0 . \end{equation*}

We first record the two facts about approximate eigenvalues that we need; both are contained in Exercise 14 of this section, but we prove them here to keep the argument self-contained.

Fact 1: every approximate eigenvalue of a bounded operator \(S\) lies in \(\operatorname{sp}(S)\). Indeed, suppose \(\beta \notin \operatorname{sp}(S)\), so \(S - \beta I\) is invertible. If \(f \in V\) with \(\|f\| = 1\), then

\begin{equation*} 1 = \|f\| = \|(S-\beta I)^{-1}(S - \beta I)f\| \leq \|(S-\beta I)^{-1}\|\,\|(S-\beta I)f\| , \end{equation*}

so \(\|(S - \beta I)f\| \geq 1/\|(S-\beta I)^{-1}\| > 0\) for every unit vector \(f\). Thus the infimum above is at least \(1/\|(S-\beta I)^{-1}\|\), which is positive, and hence \(\beta\) is not an approximate eigenvalue of \(S\).

Fact 2: every element of the spectrum of a normal operator \(S\) is an approximate eigenvalue of \(S\). Suppose \(\beta\) is not an approximate eigenvalue of \(S\). Then there exists \(c > 0\) with \(\|(S - \beta I)f\| \geq c\) for all unit vectors \(f\), and hence (by homogeneity) \(\|f\| \leq c^{-1}\|(S - \beta I)f\|\) for all \(f \in V\). The equivalence of (a) and (b) in 10.29 now shows that \(S - \beta I\) is left invertible. Furthermore \(S - \beta I\) is normal: using \((S - \beta I)^* = S^* - \overline{\beta}I\) (by 10.12) we compute

\begin{equation*} (S - \beta I)^*(S - \beta I) = S^*S - \overline{\beta}S - \beta S^* + |\beta|^2 I \quad\text{and}\quad (S - \beta I)(S - \beta I)^* = SS^* - \overline{\beta}S - \beta S^* + |\beta|^2 I , \end{equation*}

and these agree because \(S^*S = SS^*\). This is exactly the observation made in the proof of 10.56. Hence 10.55 upgrades left invertibility to invertibility, giving \(\beta \notin \operatorname{sp}(S)\).

Now suppose \(\alpha \in \operatorname{sp}(T)\). Because \(T\) is normal, Fact 2 gives a sequence \(f_1, f_2, \ldots\) in \(V\) with \(\|f_n\| = 1\) for each \(n\) and

\begin{equation*} \lim_{n \to \infty} \|(T - \alpha I)f_n\| = 0 . \end{equation*}

As noted in the proof of Fact 2, \(T - \alpha I\) is normal, so 10.53 applied to \(T - \alpha I\) gives

\begin{equation*} \|(T^* - \overline{\alpha}I)f_n\| = \|(T - \alpha I)f_n\| \longrightarrow 0 . \end{equation*}

Write \(g_n = (T - \alpha I)f_n\) and \(h_n = (T^* - \overline{\alpha}I)f_n\), so \(\|g_n\| \to 0\) and \(\|h_n\| \to 0\), and

\begin{equation*} Tf_n = \alpha f_n + g_n, \qquad T^* f_n = \overline{\alpha} f_n + h_n . \end{equation*}

Therefore

\begin{equation*} T^*Tf_n = T^*(\alpha f_n + g_n) = \alpha\,T^*f_n + T^* g_n = \alpha(\overline{\alpha}f_n + h_n) + T^*g_n = |\alpha|^2 f_n + \alpha h_n + T^* g_n . \end{equation*}

Consequently

\begin{equation*} \|(T^*T - |\alpha|^2 I)f_n\| = \|\alpha h_n + T^*g_n\| \leq |\alpha|\,\|h_n\| + \|T^*\|\,\|g_n\| \longrightarrow 0 . \end{equation*}

Because each \(f_n\) is a unit vector, this shows that \(|\alpha|^2\) is an approximate eigenvalue of \(T^*T\). By Fact 1, \(|\alpha|^2 \in \operatorname{sp}(T^*T)\), as desired.

Note that \(|\alpha|^2 \geq 0\) is a real number, so the conclusion is meaningful whether \(\mathbf{F} = \mathbf{R}\) or \(\mathbf{F} = \mathbf{C}\); also part (a) is the special case of this argument in which the approximating sequence can be taken to be constant.

Problem (10B.16)

Suppose \(\{e_k\}_{k \in \mathbf{Z}^+}\) is an orthonormal basis of a Hilbert space \(V\). Suppose also that \(T\) is a normal operator on \(V\) and \(e_k\) is an eigenvector of \(T\) for every \(k \geq 2\). Prove that \(e_1\) is an eigenvector of \(T\).

Solution

For each integer \(k \geq 2\), let \(\alpha_k \in \mathbf{F}\) be the eigenvalue of \(T\) corresponding to the eigenvector \(e_k\), so that

\begin{equation*} Te_k = \alpha_k e_k \qquad \text{for all } k \geq 2 . \end{equation*}

Because \(T\) is normal, 10.56 implies that each \(e_k\) with \(k \geq 2\) is also an eigenvector of \(T^*\), with eigenvalue \(\overline{\alpha_k}\):

\begin{equation*} T^* e_k = \overline{\alpha_k}\, e_k \qquad \text{for all } k \geq 2 . \end{equation*}

Now fix an integer \(k \geq 2\). Using the definition of the adjoint and then the equation above,

\begin{equation*} \langle Te_1, e_k \rangle = \langle e_1, T^* e_k \rangle = \langle e_1, \overline{\alpha_k} e_k \rangle = \alpha_k \langle e_1, e_k \rangle = 0 , \end{equation*}

where the last equality holds because \(\{e_k\}_{k \in \mathbf{Z}^+}\) is an orthonormal family and \(k \neq 1\).

Thus \(\langle Te_1, e_k \rangle = 0\) for every \(k \geq 2\). Because \(\{e_k\}_{k \in \mathbf{Z}^+}\) is an orthonormal basis of \(V\), the expansion 8.63(a) applied to the vector \(Te_1\) gives

\begin{equation*} Te_1 = \sum_{k=1}^{\infty} \langle Te_1, e_k \rangle e_k = \langle Te_1, e_1 \rangle e_1 , \end{equation*}

since every term with \(k \geq 2\) vanishes.

Set \(\alpha_1 = \langle Te_1, e_1 \rangle\). Then \(Te_1 = \alpha_1 e_1\) and \(\|e_1\| = 1 \neq 0\). Hence \(e_1\) is an eigenvector of \(T\), with eigenvalue \(\alpha_1\), as desired.

Two remarks. First, normality is essential: it is what allows us to move \(T\) across the inner product and land back on the eigenvectors \(e_k\), via \(T^*e_k = \overline{\alpha_k}e_k\). Without normality the conclusion can fail. For instance, let \(\{e_k\}_{k \in \mathbf{Z}^+}\) be the standard orthonormal basis of \(\ell^2\) and define a bounded operator \(T\) on \(\ell^2\) by

\begin{equation*} T(a_1, a_2, a_3, a_4, \ldots) = (0,\, a_1 + a_2,\, a_3,\, a_4, \ldots) . \end{equation*}

Then \(Te_k = e_k\) for every \(k \geq 2\), so each such \(e_k\) is an eigenvector of \(T\), but \(Te_1 = e_2\), which is not a scalar multiple of \(e_1\); thus \(e_1\) is not an eigenvector of \(T\). This \(T\) is of course not normal. Second, nothing in the argument required the eigenvalue \(\alpha_1\) to be nonzero; an eigenvector is by definition a nonzero vector, and \(e_1\) has norm \(1\), so \(e_1\) qualifies even if \(Te_1 = 0\).

Problem (10B.17)

Prove that if \(T\) is a self-adjoint operator on a Hilbert space, then \(\|T^n\| = \|T\|^n\) for every \(n \in \mathbf{Z}^+\).

Solution

Let \(V\) be the Hilbert space and let \(T \in \mathcal{B}(V)\) be self-adjoint. If \(T = 0\) the assertion is trivial (both sides equal \(0\) for every \(n \in \mathbf{Z}^+\)), so assume \(T \neq 0\); thus \(\|T\| > 0\).

Step 1: powers of \(T\) are self-adjoint. For each \(n \in \mathbf{Z}^+\), repeated use of 10.12(d) gives \((T^n)^* = (T^*)^n\), and \(T^* = T\), so \((T^n)^* = T^n\). Thus \(T^n\) is self-adjoint for every \(n \in \mathbf{Z}^+\).

Step 2: squaring doubles the exponent. Suppose \(S \in \mathcal{B}(V)\) is self-adjoint. By Exercise 6 in Section 10A (the identity \(\|S^*S\| = \|S\|^2\)),

\begin{equation*} \|S^2\| = \|S^*S\| = \|S\|^2 . \end{equation*}

In particular, applying this with \(S = T^m\) (which is self-adjoint by Step 1) gives

\begin{equation*} \|T^{2m}\| = \|T^m\|^2 \qquad \text{for every } m \in \mathbf{Z}^+ . \end{equation*}

Step 3: the result for powers of \(2\). We show by induction on \(m\) that \(\|T^{2^m}\| = \|T\|^{2^m}\) for every nonnegative integer \(m\). For \(m = 0\) this reads \(\|T\| = \|T\|\), which is true. Suppose it holds for some \(m \geq 0\). Then by Step 2 with the self-adjoint operator \(T^{2^m}\),

\begin{equation*} \|T^{2^{m+1}}\| = \|(T^{2^m})^2\| = \|T^{2^m}\|^2 = \left(\|T\|^{2^m}\right)^2 = \|T\|^{2^{m+1}} , \end{equation*}

completing the induction.

Step 4: the general case. Let \(n \in \mathbf{Z}^+\). Submultiplicativity of the operator norm (10.20) gives \(\|T^n\| \leq \|T\|^n\), so only the reverse inequality needs proof.

Choose a nonnegative integer \(m\) with \(2^m \geq n\) and set \(j = 2^m - n\). If \(j = 0\) then Step 3 immediately gives \(\|T^n\| = \|T\|^n\) and we are done. So assume \(j \geq 1\). Then, using Step 3 and submultiplicativity of the operator norm (10.20) twice,

\begin{equation*} \|T\|^{2^m} = \|T^{2^m}\| = \|T^n T^{j}\| \leq \|T^n\|\,\|T^{j}\| \leq \|T^n\|\,\|T\|^{j} = \|T^n\|\,\|T\|^{2^m - n} . \end{equation*}

Because \(\|T\| > 0\), we may divide both sides by \(\|T\|^{2^m - n}\), obtaining

\begin{equation*} \|T\|^{n} \leq \|T^n\| . \end{equation*}

Combining this with \(\|T^n\| \leq \|T\|^n\) yields \(\|T^n\| = \|T\|^n\), as desired.

Problem (10B.18)

Prove that if \(T\) is a normal operator on a Hilbert space, then \(\|T^n\| = \|T\|^n\) for every \(n \in \mathbf{Z}^+\).

Solution

Let \(V\) be the Hilbert space and let \(T \in \mathcal{B}(V)\) be normal, so \(T^*T = TT^*\). Fix \(n \in \mathbf{Z}^+\).

Step 1: \((T^*)^n T^n = (T^*T)^n\). First we check that \(T^*\) commutes with every power of \(T\). We show \(T^*T^k = T^kT^*\) by induction on \(k \in \mathbf{Z}^+\). The case \(k = 1\) is the normality hypothesis. If \(T^*T^k = T^kT^*\), then

\begin{equation*} T^*T^{k+1} = (T^*T^k)T = T^kT^*T = T^k T T^* = T^{k+1}T^* , \end{equation*}

using the induction hypothesis and then normality. This completes that induction.

Now we prove \((T^*)^n T^n = (T^*T)^n\) by induction on \(n \in \mathbf{Z}^+\). The case \(n = 1\) is a tautology. Suppose the identity holds for some \(n \in \mathbf{Z}^+\). Then

\begin{equation*} (T^*T)^{n+1} = (T^*T)^{n}(T^*T) = (T^*)^{n}T^{n}T^*T = (T^*)^{n}T^*T^{n}T = (T^*)^{n+1}T^{n+1} , \end{equation*}

where the second equality is the induction hypothesis and the third uses \(T^nT^* = T^*T^n\) from the previous paragraph. This completes the induction, so \((T^*)^n T^n = (T^*T)^n\) for all \(n \in \mathbf{Z}^+\).

Step 2: identify \((T^n)^*\). Repeated application of 10.12(d) gives \((T^n)^* = (T^*)^n\).

Step 3: compute. The operator \(T^*T\) is self-adjoint, since \((T^*T)^* = T^*(T^*)^* = T^*T\) by 10.12(d). Therefore Exercise 17 of this section applies to \(T^*T\). Using Exercise 6 in Section 10A (the identity \(\|S^*S\| = \|S\|^2\)) twice, together with Steps 1 and 2 and Exercise 17, we obtain

\begin{equation*} \begin{aligned} \|T^n\|^2 &= \|(T^n)^* T^n\| \\ &= \|(T^*)^n T^n\| \\ &= \|(T^*T)^n\| \\ &= \|T^*T\|^n \\ &= \left(\|T\|^2\right)^n \\ &= \|T\|^{2n} . \end{aligned} \end{equation*}

Taking nonnegative square roots of both ends gives \(\|T^n\| = \|T\|^n\), as desired.

Note that this contains Exercise 17 as a special case, since every self-adjoint operator is normal; but the proof above uses Exercise 17 (applied to the self-adjoint operator \(T^*T\)) as an ingredient, so the two results must be proved in this order.

Problem (10B.19)

Suppose \(T\) is an invertible operator on a Hilbert space. Prove that \(T\) is unitary if and only if \(\|T\| = \|T^{-1}\| = 1\).

Solution

Let \(V\) be the Hilbert space and let \(T \in \mathcal{B}(V)\) be invertible. Because the statement asserts that certain operator norms equal \(1\), we must assume \(V \neq \{0\}\): on the zero Hilbert space every operator equals \(0\) and has norm \(0\), yet the only operator there (namely \(0 = I\)) is unitary. So assume \(V \neq \{0\}\).

Suppose \(T\) is unitary. By the equivalence of (a) and (b) in 10.61, \(T\) is an isometry, so \(\|Tf\| = \|f\|\) for every \(f \in V\). Hence \(\|T\| = \sup\{\|Tf\| : \|f\| \leq 1\} \leq 1\), and choosing any \(f \in V\) with \(\|f\| = 1\) (possible because \(V \neq \{0\}\)) gives \(\|Tf\| = 1\), so \(\|T\| = 1\).

By the equivalence of (a) and (e) in 10.61, \(T^{-1} = T^*\), and by the equivalence of (a) and (d) in 10.61, \(T^*\) is unitary. Applying the paragraph above to the unitary operator \(T^*\) gives \(\|T^{-1}\| = \|T^*\| = 1\). Thus \(\|T\| = \|T^{-1}\| = 1\).

Suppose \(\|T\| = \|T^{-1}\| = 1\). Let \(f \in V\). Then

\begin{equation*} \|Tf\| \leq \|T\|\,\|f\| = \|f\| . \end{equation*}

On the other hand, applying the same estimate to \(T^{-1}\) and the vector \(Tf\),

\begin{equation*} \|f\| = \|T^{-1}(Tf)\| \leq \|T^{-1}\|\,\|Tf\| = \|Tf\| . \end{equation*}

Combining the two displayed inequalities gives \(\|Tf\| = \|f\|\) for every \(f \in V\), so \(T\) is an isometry.

Moreover \(T\) is invertible and hence surjective. Thus \(T\) is a surjective isometry, and the equivalence of (a) and (b) in 10.61 shows that \(T\) is unitary.

This completes both directions. Note that the invertibility hypothesis cannot be dropped from the second direction: the right shift \(T(a_1, a_2, \ldots) = (0, a_1, a_2, \ldots)\) on \(\ell^2\) satisfies \(\|T\| = 1\) and is an isometry, but it is not surjective and hence not unitary (see 10.59); of course it is also not invertible, so \(\|T^{-1}\|\) is undefined.

Problem (10B.20)

Suppose \(T\) is a bounded operator on a complex Hilbert space, with \(T = A + iB\), where \(A\) and \(B\) are self-adjoint (see 10.54). Prove that \(T\) is unitary if and only if \(T\) is normal and \(A^2 + B^2 = I\).

[If \(z = x + yi\), where \(x, y \in \mathbf{R}\), then \(|z| = 1\) if and only if \(x^2 + y^2 = 1\). Thus this exercise strengthens the analogy between the unit circle in the complex plane and the unitary operators.]

Solution

Let \(V\) be the complex Hilbert space, and let \(A, B \in \mathcal{B}(V)\) be the self-adjoint operators with \(T = A + iB\) provided by 10.54(a); recall from the proof of 10.54 that

\begin{equation*} A = \frac{T + T^*}{2}, \qquad B = \frac{T - T^*}{2i} , \end{equation*}

and that these are the unique self-adjoint operators with \(T = A + iB\).

Two computations. Since \(A^* = A\) and \(B^* = B\), parts (a) and (b) of 10.12 give

\begin{equation*} T^* = (A + iB)^* = A^* + \overline{i}\,B^* = A - iB . \end{equation*}

Therefore

\begin{equation*} \begin{aligned} T^*T &= (A - iB)(A + iB) = A^2 + iAB - iBA + B^2 = (A^2 + B^2) + i(AB - BA), \\ TT^* &= (A + iB)(A - iB) = A^2 - iAB + iBA + B^2 = (A^2 + B^2) - i(AB - BA). \end{aligned} \end{equation*}

(Here we used \((iB)(iB) = i^2B^2 = -B^2\) and \((-iB)(iB) = B^2\).)

Subtracting the two equations gives \(T^*T - TT^* = 2i(AB - BA)\), which is the identity recorded in the proof of 10.54(b); it shows once again that \(T\) is normal if and only if \(AB = BA\).

Suppose \(T\) is unitary. Then \(T^*T = TT^* = I\), so in particular \(T\) is normal (by the definition 10.50). By 10.54(b), normality of \(T\) gives \(AB = BA\), so \(AB - BA = 0\). Substituting this into the first of the two computations above,

\begin{equation*} I = T^*T = (A^2 + B^2) + i \cdot 0 = A^2 + B^2 . \end{equation*}

Thus \(T\) is normal and \(A^2 + B^2 = I\), as required.

Suppose \(T\) is normal and \(A^2 + B^2 = I\). Because \(T\) is normal, 10.54(b) gives \(AB = BA\), so \(AB - BA = 0\). Substituting into both of the computations above,

\begin{equation*} T^*T = (A^2 + B^2) + i\cdot 0 = A^2 + B^2 = I \quad\text{and}\quad TT^* = (A^2 + B^2) - i \cdot 0 = A^2 + B^2 = I . \end{equation*}

Hence \(T^*T = TT^* = I\), which by definition 10.58 says that \(T\) is unitary.

This completes both directions.

Remark on the bracketed comment. If \(A\) and \(B\) commute, the computation above shows \(T^*T = A^2 + B^2\), which is the operator analogue of \(|z|^2 = x^2 + y^2\) for \(z = x + yi\). The exercise says that being unitary (\(T^*T = TT^* = I\)) is exactly the conjunction of two conditions: the “commuting real and imaginary parts” condition, which is normality, and the “modulus one” condition \(A^2 + B^2 = I\). Note that neither condition alone suffices. For example, on \(V = \mathbf{C}\) take \(T = 2\): then \(A = 2\), \(B = 0\), so \(T\) is normal but \(A^2 + B^2 = 4I \neq I\). Conversely, on \(V = \mathbf{C}^2\) take

\begin{equation*} A = \tfrac{1}{\sqrt 2}\begin{pmatrix} 1 & 0 \\ 0 & -1\end{pmatrix}, \qquad B = \tfrac{1}{\sqrt 2}\begin{pmatrix} 0 & 1 \\ 1 & 0\end{pmatrix} . \end{equation*}

Both are self-adjoint and \(A^2 = B^2 = \tfrac12 I\), so \(A^2 + B^2 = I\). However

\begin{equation*} AB = \tfrac12\begin{pmatrix} 0 & 1 \\ -1 & 0\end{pmatrix} \qquad\text{and}\qquad BA = \tfrac12\begin{pmatrix} 0 & -1 \\ 1 & 0\end{pmatrix} , \end{equation*}

so \(AB \neq BA\). Thus \(T = A + iB\) is not normal and hence not unitary, even though \(A^2 + B^2 = I\).

Problem (10B.21)

Suppose \(T\) is a unitary operator on a complex Hilbert space such that \(T - I\) is invertible. Prove that

\begin{equation*} i(T + I)(T - I)^{-1} \end{equation*}

is a self-adjoint operator.

[The function \(z \mapsto i(z + 1)(z - 1)^{-1}\) maps \(\{z \in \mathbf{C} : |z| = 1\} \setminus \{1\}\) to \(\mathbf{R}\). Thus this exercise provides another useful illustration of the analogies showing unitary \(\approx \{z \in \mathbf{C} : |z| = 1\}\) and self-adjoint \(\approx \mathbf{R}\).]

Solution

Let \(V\) be the complex Hilbert space, let \(T \in \mathcal{B}(V)\) be unitary with \(T - I\) invertible, and set

\begin{equation*} S = i(T + I)(T - I)^{-1} . \end{equation*}

Because \(T - I\) is invertible, \((T-I)^{-1} \in \mathcal{B}(V)\), so \(S\) is a bounded operator on \(V\). We must show \(S^* = S\).

Preliminary 1: \(T^{-1} = T^*\). This is the equivalence of (a) and (e) in 10.61.

Preliminary 2: \(T + I\) and \((T - I)^{-1}\) commute. Indeed \(T + I\) and \(T - I\) commute, since

\begin{equation*} (T + I)(T - I) = T^2 - I = (T - I)(T + I) . \end{equation*}

Multiplying this equation on the left and on the right by \((T-I)^{-1}\) gives

\begin{equation*} (T - I)^{-1}(T + I) = (T + I)(T - I)^{-1} . \end{equation*}

Preliminary 3: \((T^* - I)\) is invertible, with \((T^* - I)^{-1} = \big((T - I)^{-1}\big)^*\). By 10.12, \((T - I)^* = T^* - I\). Since \(T - I\) is invertible, 10.19 says \((T-I)^*\) is invertible and \(\big((T-I)^*\big)^{-1} = \big((T-I)^{-1}\big)^*\). That is exactly the claim.

Now compute \(S^*\). Using 10.12(b) (with \(\overline{i} = -i\)), 10.12(d), and Preliminary 3,

\begin{equation*} S^* = \Big(i(T+I)(T-I)^{-1}\Big)^* = -i\,\Big((T-I)^{-1}\Big)^*\,(T+I)^* = -i\,(T^* - I)^{-1}(T^* + I) , \end{equation*}

where we also used \((T + I)^* = T^* + I\) (by 10.12).

Next we rewrite the two factors using \(T^* = T^{-1}\) (Preliminary 1). First,

\begin{equation*} T^* - I = T^{-1} - I = T^{-1}(I - T) = -T^{-1}(T - I) . \end{equation*}

Both \(T^{-1}\) and \(T - I\) are invertible, so inverting the product gives

\begin{equation*} (T^* - I)^{-1} = -\big(T^{-1}(T-I)\big)^{-1} = -(T - I)^{-1}T . \end{equation*}

Second,

\begin{equation*} T^* + I = T^{-1} + I = T^{-1}(I + T) = T^{-1}(T + I) . \end{equation*}

Substituting these two expressions into the formula for \(S^*\),

\begin{equation*} \begin{aligned} S^* &= -i\,\Big(-(T-I)^{-1}T\Big)\Big(T^{-1}(T+I)\Big) \\ &= i\,(T-I)^{-1}\,T\,T^{-1}\,(T+I) \\ &= i\,(T-I)^{-1}(T+I) . \end{aligned} \end{equation*}

By Preliminary 2, \((T-I)^{-1}(T+I) = (T+I)(T-I)^{-1}\), and therefore

\begin{equation*} S^* = i(T + I)(T - I)^{-1} = S . \end{equation*}

Hence \(S = i(T+I)(T-I)^{-1}\) is self-adjoint, as desired.

Remark. The bracketed comment in the statement is the scalar shadow of this computation. If \(|z| = 1\) and \(z \neq 1\), then \(\overline{z} = 1/z\), so

\begin{equation*} \overline{\,i\,\frac{z+1}{z-1}\,} = -i\,\frac{\overline{z}+1}{\overline{z}-1} = -i\,\frac{\frac1z + 1}{\frac1z - 1} = -i\,\frac{1 + z}{1 - z} = i\,\frac{z+1}{z-1} , \end{equation*}

so \(i(z+1)(z-1)^{-1}\) is its own conjugate and hence real. The proof above is the same manipulation with \(z\) replaced by \(T\), complex conjugation replaced by taking adjoints, and the only extra care being that operators need not commute, which is handled by Preliminary 2. Combining this exercise with 10.49 and 10.62, one sees the map \(T \mapsto i(T+I)(T-I)^{-1}\) (the Cayley transform) carrying unitary operators without \(1\) in their spectrum to self-adjoint operators, mirroring the Mobius transformation carrying the punctured unit circle to \(\mathbf{R}\).

Problem (10B.22)

Suppose \(T\) is a self-adjoint operator on a complex Hilbert space. Prove that

\begin{equation*} (T + iI)(T - iI)^{-1} \end{equation*}

is a unitary operator.

[The function \(z \mapsto (z + i)(z - i)^{-1}\) maps \(\mathbb{R}\) to \(\{z \in \mathbb{C} : |z| = 1\} \setminus \{1\}\). Thus this exercise provides another useful illustration of the analogies showing (a) unitary \(\Longleftrightarrow\) \(\{z \in \mathbb{C} : |z| = 1\}\); (b) self-adjoint \(\Longleftrightarrow\) \(\mathbb{R}\).]

Solution

Let \(V\) be the complex Hilbert space and let \(T\) be self-adjoint; thus \(T \in \mathcal{B}(V)\) and \(T^* = T\) (self-adjointness is defined only for bounded operators, because the adjoint is).

Step 1: the expression makes sense. Because \(T\) is self-adjoint, 10.49 gives \(\operatorname{sp}(T) \subseteq \mathbb{R}\). Hence \(i \notin \operatorname{sp}(T)\) and \(-i \notin \operatorname{sp}(T)\), which means that both \(T - iI\) and \(T + iI\) are invertible elements of \(\mathcal{B}(V)\). In particular \((T - iI)^{-1} \in \mathcal{B}(V)\), so

\begin{equation*} U := (T + iI)(T - iI)^{-1} \end{equation*}

is a well-defined bounded operator on \(V\).

Step 2: computing \(U^*\). By 10.12(b) and 10.12(c), \((\alpha I)^* = \overline{\alpha} I\) for every \(\alpha \in \mathbb{C}\). Because the adjoint is additive [10.12(a)] and \(T^* = T\),

\begin{equation*} (T - iI)^* = T^* + iI = T + iI, \qquad (T + iI)^* = T^* - iI = T - iI . \end{equation*}

Also, if \(S \in \mathcal{B}(V)\) is invertible then \((S^{-1})^* = (S^*)^{-1}\): applying the adjoint to \(S S^{-1} = S^{-1} S = I\) and using 10.12(c) and 10.12(d) gives \((S^{-1})^* S^* = S^* (S^{-1})^* = I\), so \(S^*\) is invertible with inverse \((S^{-1})^*\). Therefore

\begin{equation*} U^* = \bigl((T - iI)^{-1}\bigr)^* (T + iI)^* = \bigl((T - iI)^*\bigr)^{-1}(T - iI) = (T + iI)^{-1}(T - iI). \end{equation*}

Step 3: everything in sight commutes. Expanding,

\begin{equation*} (T + iI)(T - iI) = T^2 - iT + iT + I = T^2 + I = (T - iI)(T + iI). \end{equation*}

Thus \(A := T + iI\) and \(B := T - iI\) commute. If two operators commute and one of them is invertible, then that inverse also commutes with the other: from \(AB = BA\) we get \(B^{-1}(AB)B^{-1} = B^{-1}(BA)B^{-1}\), i.e. \(B^{-1}A = A B^{-1}\). Applying this twice shows that \(A^{-1}\) and \(B^{-1}\) commute with each other and with \(A\) and \(B\).

Step 4: \(U\) is unitary. Using the commutation relations from Step 3,

\begin{equation*} U^* U = A^{-1} B \, A \, B^{-1} = A^{-1} A \, B \, B^{-1} = I, \end{equation*}

and

\begin{equation*} U U^* = A B^{-1} \, A^{-1} B = A A^{-1} \, B^{-1} B = I . \end{equation*}

Hence \(U^* U = U U^* = I\), so \(U\) is unitary by the definition 10.58.

Remark: an alternative route. One can also argue metrically. For \(f \in V\) we have \(\langle Tf, f \rangle \in \mathbb{R}\) by 10.48, and \(\langle f, Tf \rangle = \overline{\langle Tf, f \rangle} = \langle Tf, f \rangle\). Hence

\begin{equation*} \|(T + iI)f\|^2 = \|Tf\|^2 + \langle Tf, if \rangle + \langle if, Tf \rangle + \|f\|^2 = \|Tf\|^2 - i\langle Tf, f\rangle + i \langle f, Tf \rangle + \|f\|^2 = \|Tf\|^2 + \|f\|^2, \end{equation*}

and the same computation with \(-i\) in place of \(i\) gives \(\|(T - iI)f\|^2 = \|Tf\|^2 + \|f\|^2\). Thus \(\|(T + iI)f\| = \|(T - iI)f\|\) for every \(f \in V\). Replacing \(f\) by \((T - iI)^{-1}g\) shows that \(\|Ug\| = \|g\|\) for all \(g \in V\), so \(U\) is an isometry; and \(U\) is surjective because it is the composition of the two bijections \((T - iI)^{-1}\) and \(T + iI\) of \(V\) onto \(V\). A surjective isometry is unitary by 10.61.

Problem (10B.23)

For \(T\) a bounded operator on a Banach space, define \(e^T\) by

\begin{equation*} e^T = \sum_{k=0}^{\infty} \frac{T^k}{k!} . \end{equation*}

(a) Prove that if \(T\) is a bounded operator on a Banach space \(V\), then the infinite sum above converges in \(\mathcal{B}(V)\) and \(\|e^T\| \le e^{\|T\|}\).

(b) Prove that if \(S, T\) are bounded operators on a Banach space \(V\) such that \(ST = TS\), then \(e^S e^T = e^{S+T}\).

(c) Prove that if \(T\) is a self-adjoint operator on a complex Hilbert space, then \(e^{iT}\) is unitary.

Solution

Throughout, \(T^0 = I\) by convention, and we use repeatedly that \(\|AB\| \le \|A\|\,\|B\|\) for \(A, B \in \mathcal{B}(V)\), which holds because \(\|ABf\| \le \|A\|\,\|Bf\| \le \|A\|\,\|B\|\,\|f\|\) for every \(f \in V\). Iterating gives \(\|A^k\| \le \|A\|^k\) for every \(k \in \mathbb{Z}^+\); and \(\|A^0\| = \|I\| \le 1 = \|A\|^0\) as well (with \(\|I\| = 1\) unless \(V = \{0\}\), in which case \(\|I\| = 0\)). So

\begin{equation*} \|T^k\| \le \|T\|^k \quad \text{for every integer } k \ge 0 . \end{equation*}

(a) Let \(V\) be a Banach space and \(T \in \mathcal{B}(V)\). By 6.47 (applied with \(W = V\)), \(\mathcal{B}(V)\) is a Banach space. Now

\begin{equation*} \sum_{k=0}^{\infty} \left\| \frac{T^k}{k!} \right\| = \sum_{k=0}^{\infty} \frac{\|T^k\|}{k!} \le \sum_{k=0}^{\infty} \frac{\|T\|^k}{k!} = e^{\|T\|} < \infty . \end{equation*}

Thus the series \(\sum_{k=0}^{\infty} T^k / k!\) is absolutely convergent in the Banach space \(\mathcal{B}(V)\), and hence it converges, by 6.41. So \(e^T\) is a well-defined element of \(\mathcal{B}(V)\).

For the norm estimate, let \(A_n = \sum_{k=0}^{n} T^k/k!\) denote the \(n\)-th partial sum. Then

\begin{equation*} \|A_n\| \le \sum_{k=0}^{n} \frac{\|T^k\|}{k!} \le \sum_{k=0}^{n} \frac{\|T\|^k}{k!} \le e^{\|T\|} \end{equation*}

for every \(n\). Because \(A_n \to e^T\) in \(\mathcal{B}(V)\) and the norm is a continuous function (by the reverse triangle inequality), we conclude that \(\|e^T\| = \lim_{n \to \infty} \|A_n\| \le e^{\|T\|}\).

(b) Suppose \(S, T \in \mathcal{B}(V)\) and \(ST = TS\). Because \(S\) and \(T\) commute, the binomial theorem holds for them:

\begin{equation*} (S + T)^m = \sum_{k=0}^{m} \binom{m}{k} S^k T^{m-k} \quad \text{for every integer } m \ge 0 . \end{equation*}

This is proved by induction on \(m\) exactly as for numbers: the case \(m = 0\) is trivial, and the induction step uses only that \(T^j S = S T^j\) (which follows from \(ST = TS\) by induction on \(j\)) to move factors of \(S\) past factors of \(T\), together with Pascal’s identity.

Set

\begin{equation*} A_n = \sum_{j=0}^{n} \frac{S^j}{j!}, \qquad B_n = \sum_{k=0}^{n} \frac{T^k}{k!}, \qquad C_n = \sum_{m=0}^{n} \frac{(S+T)^m}{m!} . \end{equation*}

By part (a), \(A_n \to e^S\), \(B_n \to e^T\), and \(C_n \to e^{S+T}\) in \(\mathcal{B}(V)\).

Multiplying out and using the binomial theorem,

\begin{equation*} A_n B_n = \sum_{j=0}^{n} \sum_{k=0}^{n} \frac{S^j T^k}{j!\,k!}, \qquad C_n = \sum_{m=0}^{n} \frac{1}{m!}\sum_{k=0}^{m} \binom{m}{k} S^k T^{m-k} = \sum_{\substack{j, k \ge 0 \\ j + k \le n}} \frac{S^j T^k}{j!\,k!} . \end{equation*}

Every index pair occurring in \(C_n\) also occurs in \(A_n B_n\) (since \(j + k \le n\) forces \(j \le n\) and \(k \le n\)). Hence

\begin{equation*} A_n B_n - C_n = \sum_{(j,k) \in D_n} \frac{S^j T^k}{j!\,k!}, \qquad \text{where } D_n = \{(j,k) : 0 \le j, k \le n \text{ and } j + k > n\} . \end{equation*}

Writing \(a = \|S\|\) and \(b = \|T\|\), the triangle inequality and \(\|S^j T^k\| \le a^j b^k\) give

\begin{equation*} \|A_n B_n - C_n\| \le \sum_{(j,k) \in D_n} \frac{a^j b^k}{j!\,k!} = \left(\sum_{j=0}^{n} \frac{a^j}{j!}\right)\left(\sum_{k=0}^{n} \frac{b^k}{k!}\right) - \sum_{\substack{j,k \ge 0 \\ j+k \le n}} \frac{a^j b^k}{j!\,k!} , \end{equation*}

where the last equality is the same index bookkeeping as above, now for nonnegative real numbers. As \(n \to \infty\) the first product tends to \(e^a e^b\), while the subtracted sum equals \(\sum_{m=0}^{n} (a+b)^m / m!\) by the ordinary binomial theorem and hence tends to \(e^{a+b} = e^a e^b\). Therefore \(\|A_n B_n - C_n\| \to 0\).

Finally, \(A_n B_n \to e^S e^T\): indeed

\begin{equation*} \|A_n B_n - e^S e^T\| \le \|A_n - e^S\|\,\|B_n\| + \|e^S\|\,\|B_n - e^T\| \to 0, \end{equation*}

because \(\|B_n\| \le e^b\) for all \(n\) by the partial-sum estimate in part (a). Combining,

\begin{equation*} e^S e^T = \lim_{n \to \infty} A_n B_n = \lim_{n \to \infty} C_n = e^{S+T} . \end{equation*}

(c) Let \(V\) be a complex Hilbert space and \(T \in \mathcal{B}(V)\) with \(T^* = T\). First we record that \((e^A)^* = e^{A^*}\) for every \(A \in \mathcal{B}(V)\). Indeed, the adjoint map \(A \mapsto A^*\) satisfies \((A + B)^* = A^* + B^*\), \((\alpha A)^* = \overline{\alpha} A^*\), \(I^* = I\), and \((AB)^* = B^* A^*\) (all by 10.12), and \(\|A^*\| = \|A\|\) (by 10.11); since \(A \mapsto A^*\) is additive and norm-preserving, \(\|A^* - B^*\| = \|(A-B)^*\| = \|A - B\|\), so \(A \mapsto A^*\) is continuous on \(\mathcal{B}(V)\). Since \((A^k)^* = (A^*)^k\) for every integer \(k \ge 0\) (immediate from 10.12(c) and 10.12(d) by induction) and \(1/k!\) is real,

\begin{equation*} \left( \sum_{k=0}^{n} \frac{A^k}{k!} \right)^{\!*} = \sum_{k=0}^{n} \frac{(A^*)^k}{k!} , \end{equation*}

and letting \(n \to \infty\) (using continuity of the adjoint) gives \((e^A)^* = e^{A^*}\).

Apply this with \(A = iT\). Since \((iT)^* = \overline{i}\, T^* = -iT\), we get

\begin{equation*} (e^{iT})^* = e^{-iT} . \end{equation*}

Moreover \((iT)(-iT) = T^2 = (-iT)(iT)\), so \(iT\) and \(-iT\) commute. Hence part (b) applies and yields

\begin{equation*} (e^{iT})^* e^{iT} = e^{-iT} e^{iT} = e^{-iT + iT} = e^{0} = I, \end{equation*}

\begin{equation*} e^{iT} (e^{iT})^* = e^{iT} e^{-iT} = e^{0} = I, \end{equation*}

where \(e^0 = I\) because all terms of the series with \(k \ge 1\) vanish and the \(k = 0\) term is \(I\). Thus \(e^{iT}\) is unitary by definition 10.58.

Problem (10B.24)

A bounded operator \(T\) on a Hilbert space is called a partial isometry if

\begin{equation*} \|Tf\| = \|f\| \ \text{ for all } f \in (\operatorname{null} T)^{\perp}. \end{equation*}

Suppose \((X, \mathcal{S}, \mu)\) is a \(\sigma\)-finite measure space and \(h \in L^{\infty}(\mu)\). As usual, let \(M_h \in \mathcal{B}\bigl(L^2(\mu)\bigr)\) denote the multiplication operator defined by \(M_h f = fh\). Prove that \(M_h\) is a partial isometry if and only if there exists a set \(E \in \mathcal{S}\) such that \(|h| = \chi_E\).

Solution

Fix an \(\mathcal{S}\)-measurable representative of \(h\) with \(|h(x)| \le \|h\|_{\infty}\) for all \(x \in X\); all statements about \(h\) below are understood up to sets of \(\mu\)-measure \(0\), which is exactly the sense in which the asserted identity \(|h| = \chi_E\) is to be read (equality in \(L^{\infty}(\mu)\)). Recall that \(M_h\) is bounded, with \(\|M_h\| \le \|h\|_{\infty}\).

Put

\begin{equation*} Z = \{x \in X : h(x) = 0\}, \qquad W = X \setminus Z = \{x \in X : h(x) \ne 0\} ; \end{equation*}

both belong to \(\mathcal{S}\) because \(h\) is \(\mathcal{S}\)-measurable.

Step 1: identification of \(\operatorname{null} M_h\) and its orthogonal complement. For \(f \in L^2(\mu)\) we have \(M_h f = 0\) if and only if \(fh = 0\) almost everywhere, which happens if and only if \(f = 0\) almost everywhere on \(W\). Thus

\begin{equation*} \operatorname{null} M_h = \{ f \in L^2(\mu) : f = 0 \text{ a.e. on } W \} . \end{equation*}

We claim that

\begin{equation*} (\operatorname{null} M_h)^{\perp} = \{ f \in L^2(\mu) : f = 0 \text{ a.e. on } Z \} . \end{equation*}

For the inclusion \(\supseteq\): if \(f = 0\) a.e. on \(Z\) and \(g \in \operatorname{null} M_h\), then \(g = 0\) a.e. on \(W\), so \(f \overline{g} = 0\) almost everywhere on \(X\) and hence \(\langle f, g \rangle = \int_X f \overline{g} \, d\mu = 0\). For the inclusion \(\subseteq\): suppose \(f \in (\operatorname{null} M_h)^{\perp}\). The function \(g = \chi_Z f\) lies in \(L^2(\mu)\) and vanishes on \(W\), so \(g \in \operatorname{null} M_h\); hence

\begin{equation*} 0 = \langle f, g \rangle = \int_X f \, \overline{\chi_Z f} \, d\mu = \int_Z |f|^2 \, d\mu , \end{equation*}

which forces \(f = 0\) almost everywhere on \(Z\). This proves the claim.

Because \(\|M_h f\|^2 = \int_X |fh|^2 \, d\mu = \int_X |f|^2 |h|^2 \, d\mu\) and \(\|f\|^2 = \int_X |f|^2 \, d\mu\), we conclude that \(M_h\) is a partial isometry if and only if

\begin{equation*} \int_X |f|^2 |h|^2 \, d\mu = \int_X |f|^2 \, d\mu \quad \text{for every } f \in L^2(\mu) \text{ with } f = 0 \text{ a.e. on } Z , \end{equation*}

that is, if and only if

\begin{equation*} (\ast) \qquad \int_W |f|^2 \bigl(|h|^2 - 1\bigr) \, d\mu = 0 \quad \text{for every such } f . \end{equation*}

Step 2: suppose \(|h| = \chi_E\) for some \(E \in \mathcal{S}\). Then \(|h(x)| = 1\) for almost every \(x \in E\) and \(h(x) = 0\) for almost every \(x \notin E\). Hence \(W = E\) up to a set of measure \(0\), and \(|h|^2 - 1 = 0\) almost everywhere on \(W\). So the integrand in \((\ast)\) vanishes almost everywhere and \((\ast)\) holds. Therefore \(M_h\) is a partial isometry.

Step 3: conversely, suppose \(M_h\) is a partial isometry; we show \(|h| = 1\) almost everywhere on \(W\). Suppose not. Then \(\mu\bigl(\{x \in W : |h(x)| \ne 1\}\bigr) > 0\), and since

\begin{equation*} \{x \in W : |h(x)| \ne 1\} = \bigcup_{n=1}^{\infty} \Bigl( \bigl\{ x \in W : |h(x)| \ge 1 + \tfrac1n \bigr\} \cup \bigl\{ x \in W : 0 < |h(x)| \le 1 - \tfrac1n \bigr\} \Bigr), \end{equation*}

countable subadditivity gives an \(n \in \mathbb{Z}^+\) and a set \(F \in \mathcal{S}\) with \(F \subseteq W\), \(\mu(F) > 0\), and either \(|h| \ge 1 + \frac1n\) on \(F\) or \(|h| \le 1 - \frac1n\) on \(F\). Because \(\mu\) is \(\sigma\)-finite, write \(X = \bigcup_{j=1}^{\infty} X_j\) with \(\mu(X_j) < \infty\); since \(\mu(F) > 0\), there is a \(j\) with \(\mu(F \cap X_j) > 0\). Set \(A = F \cap X_j\), so

\begin{equation*} A \in \mathcal{S}, \qquad A \subseteq W, \qquad 0 < \mu(A) < \infty . \end{equation*}

Then \(f = \chi_A \in L^2(\mu)\) and \(f = 0\) on \(Z\), so \(f \in (\operatorname{null} M_h)^{\perp}\), while \(\|f\|^2 = \mu(A)\). In the first case,

\begin{equation*} \|M_h f\|^2 = \int_A |h|^2 \, d\mu \ge \Bigl(1 + \tfrac{1}{n}\Bigr)^2 \mu(A) > \mu(A) = \|f\|^2 , \end{equation*}

and in the second case,

\begin{equation*} \|M_h f\|^2 = \int_A |h|^2 \, d\mu \le \Bigl(1 - \tfrac{1}{n}\Bigr)^2 \mu(A) < \mu(A) = \|f\|^2 , \end{equation*}

the strict inequalities holding because \(0 < \mu(A) < \infty\) and \(0 \le (1 - \frac{1}{n})^2 < 1 < (1 + \frac{1}{n})^2\). Either way we contradict \(\|M_h f\| = \|f\|\).

Hence \(|h| = 1\) almost everywhere on \(W\), while \(|h| = 0\) on \(Z\) by the definition of \(Z\). Therefore \(|h| = \chi_W\) almost everywhere, and \(E = W \in \mathcal{S}\) is the desired set.

Problem (10B.25)

Suppose \(T\) is an isometry on a Hilbert space. Prove that \(T^*\) is a partial isometry.

Solution

Let \(V\) be the Hilbert space and let \(T \in \mathcal{B}(V)\) be an isometry, so \(\|Tf\| = \|f\|\) for every \(f \in V\). Recall that \(T^*\) is a partial isometry means \(\|T^* g\| = \|g\|\) for every \(g \in (\operatorname{null} T^*)^{\perp}\).

Step 1: \(T^* T = I\). This is the equivalence of (a) and (c) in 10.60.

Step 2: \(\operatorname{range} T\) is closed. Suppose \(g \in \overline{\operatorname{range} T}\), say \(T f_n \to g\) with \(f_1, f_2, \ldots \in V\). Because \(T\) is an isometry and linear,

\begin{equation*} \|f_n - f_m\| = \|T(f_n - f_m)\| = \|Tf_n - Tf_m\| , \end{equation*}

and the right side tends to \(0\) as \(n, m \to \infty\) because the convergent sequence \(Tf_1, Tf_2, \ldots\) is Cauchy. Hence \(f_1, f_2, \ldots\) is a Cauchy sequence in \(V\); since \(V\) is complete, \(f_n \to f\) for some \(f \in V\). Because \(T\) is continuous, \(Tf_n \to Tf\), and therefore \(g = Tf \in \operatorname{range} T\). Thus \(\operatorname{range} T\) is closed.

Step 3: \((\operatorname{null} T^*)^{\perp} = \operatorname{range} T\). By 10.13(d) we have \((\operatorname{null} T^*)^{\perp} = \overline{\operatorname{range} T}\), which equals \(\operatorname{range} T\) by Step 2. (Equivalently: \(\operatorname{null} T^* = (\operatorname{range} T)^{\perp}\) by 10.13(a), and taking orthogonal complements of the closed subspace \(\operatorname{range} T\) gives the same conclusion via 8.41.)

Step 4: \(T^*\) is isometric on \((\operatorname{null} T^*)^{\perp}\). Let \(g \in (\operatorname{null} T^*)^{\perp}\). By Step 3 there exists \(f \in V\) with \(g = Tf\). Then by Step 1,

\begin{equation*} T^* g = T^* T f = f , \end{equation*}

and hence, using that \(T\) is an isometry,

\begin{equation*} \|T^* g\| = \|f\| = \|Tf\| = \|g\| . \end{equation*}

Thus \(\|T^* g\| = \|g\|\) for every \(g \in (\operatorname{null} T^*)^{\perp}\), which says exactly that \(T^*\) is a partial isometry.

Problem (10B.26)

Suppose \(T\) is a bounded operator on a Hilbert space \(V\). Prove that \(T\) is a partial isometry if and only if \(T^* T = P_U\) for some closed subspace \(U\) of \(V\).

Solution

Here \(P_U\) denotes the orthogonal projection of \(V\) onto the closed subspace \(U\) (see Section 8B). We use freely the following facts: \(P_U\) is self-adjoint (last bullet point of 10.45); \(\operatorname{range} P_U = U\) and \(\operatorname{null} P_U = U^{\perp}\) (8.45(a)); every \(f \in V\) decomposes uniquely as \(f = P_U f + (f - P_U f)\) with \(P_U f \in U\) and \(f - P_U f \in U^{\perp}\) (8.43); and \(\overline{W} = (W^{\perp})^{\perp}\) for every subspace \(W\) (8.41).

A preliminary computation. If \(W\) is a closed subspace of \(V\) and \(f \in V\), then \(f = P_W f + (f - P_W f)\) with the two summands orthogonal, so

\begin{equation*} \langle f, P_W f \rangle = \langle P_W f, P_W f \rangle + \langle f - P_W f, P_W f\rangle = \|P_W f\|^2 . \end{equation*}

Also, for any \(S \in \mathcal{B}(V)\) and \(f \in V\), the definition of the adjoint gives

\begin{equation*} \|Sf\|^2 = \langle Sf, Sf \rangle = \langle f, S^* S f \rangle . \end{equation*}

First suppose \(T^* T = P_U\) for some closed subspace \(U\) of \(V\). Then for every \(f \in V\),

\begin{equation*} \|Tf\|^2 = \langle f, T^* T f \rangle = \langle f, P_U f \rangle = \|P_U f\|^2 , \end{equation*}

so

\begin{equation*} (\dagger) \qquad \|Tf\| = \|P_U f\| \quad \text{for all } f \in V . \end{equation*}

In particular \(Tf = 0\) if and only if \(P_U f = 0\), so by 8.45(a),

\begin{equation*} \operatorname{null} T = \operatorname{null} P_U = U^{\perp} . \end{equation*}

Taking orthogonal complements and using 8.41 together with the fact that \(U\) is closed,

\begin{equation*} (\operatorname{null} T)^{\perp} = (U^{\perp})^{\perp} = \overline{U} = U . \end{equation*}

Now if \(f \in (\operatorname{null} T)^{\perp} = U\), then \(P_U f = f\), so \((\dagger)\) gives \(\|Tf\| = \|f\|\). Hence \(T\) is a partial isometry.

Conversely, suppose \(T\) is a partial isometry. Because \(T\) is bounded, \(\operatorname{null} T\) is a closed subspace of \(V\). Set

\begin{equation*} U = (\operatorname{null} T)^{\perp} , \end{equation*}

which is a closed subspace of \(V\) (orthogonal complements are always closed subspaces, by 8.40(a)). Note that, because \(\operatorname{null} T\) is closed, 8.41 gives

\begin{equation*} U^{\perp} = \bigl((\operatorname{null} T)^{\perp}\bigr)^{\perp} = \overline{\operatorname{null} T} = \operatorname{null} T . \end{equation*}

Let \(f \in V\). By 8.43, \(f = g + h\) with \(g = P_U f \in U\) and \(h = f - P_U f \in U^{\perp} = \operatorname{null} T\). Hence \(Th = 0\) and so \(Tf = Tg\). Since \(g \in U = (\operatorname{null} T)^{\perp}\) and \(T\) is a partial isometry, \(\|Tg\| = \|g\|\). Therefore, using the preliminary computation twice,

\begin{equation*} \langle f, T^* T f \rangle = \|Tf\|^2 = \|Tg\|^2 = \|g\|^2 = \|P_U f\|^2 = \langle f, P_U f \rangle . \end{equation*}

Thus

\begin{equation*} \langle f, (T^* T - P_U) f \rangle = 0 \quad \text{for every } f \in V . \end{equation*}

The operator \(S = T^* T - P_U\) is self-adjoint, because \((T^* T)^* = T^* (T^*)^* = T^* T\) and \((P_U)^* = P_U\). For a self-adjoint \(S\) we have \(\langle Sf, f \rangle = \overline{\langle f, Sf \rangle} = 0\) for all \(f \in V\), and hence \(S = 0\) by 10.46 (part (a) if \(\mathbb{F} = \mathbb{C}\), part (b) if \(\mathbb{F} = \mathbb{R}\), the latter applying because \(S\) is self-adjoint).

Therefore \(T^* T = P_U\) with \(U = (\operatorname{null} T)^{\perp}\) a closed subspace of \(V\), completing the proof.

Exercises 10C

Problem (10C.1)

Prove that if \(T\) is a compact operator on a Hilbert space \(V\) and \(e_1, e_2, \ldots\) is an orthonormal sequence in \(V\), then \(\lim_{n \to \infty} T e_n = 0\).

Solution

Throughout, \(\langle \cdot, \cdot \rangle\) denotes the inner product on \(V\) and \(\|\cdot\|\) the associated norm.

Step 1: an orthonormal sequence tends weakly to \(0\).

Let \(h \in V\). Applying Bessel’s inequality (8.57) to the orthonormal family \(\{e_n\}_{n \in \mathbf{Z}^+}\) and the vector \(h\) gives

\begin{equation*} \sum_{n=1}^{\infty} |\langle h, e_n \rangle|^2 \le \|h\|^2 < \infty . \end{equation*}

A convergent series of nonnegative numbers has terms tending to \(0\), so

\begin{equation*} \lim_{n \to \infty} \langle h, e_n \rangle = 0 , \qquad \text{and hence} \qquad \lim_{n \to \infty} \langle e_n, h \rangle = 0 \end{equation*}

for every \(h \in V\), because \(\langle e_n, h \rangle = \overline{\langle h, e_n \rangle}\).

Step 2: the compactness argument.

Suppose, seeking a contradiction, that \(T e_n\) does not converge to \(0\). Then there exist \(\varepsilon > 0\) and a subsequence \(e_{n_1}, e_{n_2}, \ldots\) such that

\begin{equation*} \|T e_{n_k}\| \ge \varepsilon \qquad \text{for every } k \in \mathbf{Z}^+ . \end{equation*}

The sequence \(e_{n_1}, e_{n_2}, \ldots\) is bounded, because \(\|e_{n_k}\| = 1\) for every \(k\). Since \(T\) is compact, the definition of compactness (10.66) provides a further subsequence \(e_{n_{k_1}}, e_{n_{k_2}}, \ldots\) and a vector \(u \in V\) with

\begin{equation*} \lim_{j \to \infty} T e_{n_{k_j}} = u . \end{equation*}

Because the inner product is continuous in its first slot (which follows from the Cauchy–Schwarz inequality: \(|\langle v, u\rangle - \langle v^{\prime}, u \rangle| \le \|v - v^{\prime}\| \, \|u\|\)), we get

\begin{equation*} \lim_{j \to \infty} \langle T e_{n_{k_j}}, u \rangle = \langle u, u \rangle = \|u\|^2 . \end{equation*}

On the other hand, \(T\) is bounded (10.68), so its adjoint \(T^*\) exists, and for every \(n\) we have \(\langle T e_n, u \rangle = \langle e_n, T^* u \rangle\). Taking \(h = T^* u\) in Step 1 shows that

\begin{equation*} \lim_{n \to \infty} \langle T e_n, u \rangle = \lim_{n \to \infty} \langle e_n, T^* u \rangle = 0 , \end{equation*}

and in particular the limit along the subsequence is \(0\). Comparing the two displayed limits gives \(\|u\|^2 = 0\), so \(u = 0\).

But then \(\lim_{j \to \infty} \|T e_{n_{k_j}}\| = \|u\| = 0\), which contradicts \(\|T e_{n_{k_j}}\| \ge \varepsilon > 0\) for all \(j\).

This contradiction shows that no such \(\varepsilon\) and subsequence exist. Hence \(\lim_{n \to \infty} T e_n = 0\), as desired.

Problem (10C.2)

Prove that if \(T\) is a compact operator on \(L^2([0,1])\), then

\begin{equation*} \lim_{n \to \infty} \sqrt{n} \, \|T(x^n)\|_2 = 0 , \end{equation*}

where \(x^n\) means the element of \(L^2([0,1])\) defined by \(x \mapsto x^n\).

Solution

Step 1: normalize.

For \(n \in \mathbf{Z}^+\) we have

\begin{equation*} \|x^n\|_2^2 = \int_0^1 x^{2n} \, dx = \frac{1}{2n+1} , \end{equation*}

so \(\|x^n\|_2 = (2n+1)^{-1/2}\). Define

\begin{equation*} g_n = \sqrt{2n+1} \, x^n \in L^2([0,1]) , \end{equation*}

so that \(\|g_n\|_2 = 1\) for every \(n\).

Step 2: \(\lim_{n \to \infty} \langle g_n, h \rangle = 0\) for every \(h \in L^2([0,1])\).

Fix \(h \in L^2([0,1])\) and let \(\varepsilon > 0\).

Because \(|h|^2\) is integrable on \([0,1]\), the Dominated Convergence Theorem (3.31) applied to the functions \(|h|^2 \chi_{[1 - 1/m, \, 1]}\) (which converge pointwise to \(0\) at every \(x < 1\), hence almost everywhere, and are dominated by \(|h|^2\)) shows that

\begin{equation*} \lim_{\delta \downarrow 0} \int_{1-\delta}^{1} |h|^2 = 0 . \end{equation*}

Choose \(\delta \in (0,1)\) such that

\begin{equation*} \Bigl( \int_{1-\delta}^{1} |h|^2 \Bigr)^{1/2} < \tfrac{\varepsilon}{2} . \end{equation*}

Now split the inner product:

\begin{equation*} |\langle g_n, h \rangle| \le \sqrt{2n+1} \int_0^{1-\delta} x^n |h| \, dx \; + \; \sqrt{2n+1} \int_{1-\delta}^{1} x^n |h| \, dx . \end{equation*}

For the second term, the Cauchy–Schwarz inequality gives

\begin{equation*} \int_{1-\delta}^{1} x^n |h| \, dx \le \Bigl( \int_{1-\delta}^{1} x^{2n} \Bigr)^{1/2} \Bigl( \int_{1-\delta}^{1} |h|^2 \Bigr)^{1/2} \le \frac{1}{\sqrt{2n+1}} \cdot \frac{\varepsilon}{2} , \end{equation*}

so the second term is less than \(\varepsilon/2\) for every \(n\).

For the first term, \(x^n \le (1-\delta)^n\) on \([0, 1-\delta]\), so by the Cauchy–Schwarz inequality on the probability space \([0,1]\),

\begin{equation*} \int_0^{1-\delta} x^n |h| \, dx \le (1-\delta)^n \int_0^1 |h| \le (1-\delta)^n \|h\|_2 . \end{equation*}

Since \(0 < 1 - \delta < 1\), we have \(\lim_{n \to \infty} \sqrt{2n+1} \, (1-\delta)^n = 0\) (an exponentially decaying factor beats the growth of \(\sqrt{2n+1}\)). Hence there exists \(N \in \mathbf{Z}^+\) such that \(\sqrt{2n+1} \, (1-\delta)^n \|h\|_2 < \varepsilon/2\) for all \(n \ge N\).

Combining, \(|\langle g_n, h \rangle| < \varepsilon\) for all \(n \ge N\). As \(\varepsilon > 0\) was arbitrary, \(\lim_{n \to \infty} \langle g_n, h \rangle = 0\).

Step 3: compactness forces \(\|T g_n\|_2 \to 0\).

Suppose not. Then there exist \(\varepsilon_0 > 0\) and a subsequence with \(\|T g_{n_k}\|_2 \ge \varepsilon_0\) for all \(k\). The sequence \(g_{n_1}, g_{n_2}, \ldots\) is bounded (each has norm \(1\)), so compactness of \(T\) (10.66) yields a further subsequence and a \(u \in L^2([0,1])\) with \(\lim_{j \to \infty} T g_{n_{k_j}} = u\). Continuity of the inner product then gives

\begin{equation*} \lim_{j \to \infty} \langle T g_{n_{k_j}}, u \rangle = \|u\|_2^2 . \end{equation*}

But \(\langle T g_n, u \rangle = \langle g_n, T^* u \rangle\), which tends to \(0\) by Step 2 applied to \(h = T^* u\). Hence \(\|u\|_2 = 0\), so \(u = 0\) and \(\|T g_{n_{k_j}}\|_2 \to 0\), contradicting \(\|T g_{n_{k_j}}\|_2 \ge \varepsilon_0\). Therefore

\begin{equation*} \lim_{n \to \infty} \|T g_n\|_2 = 0 . \end{equation*}

(This step is exactly Exercise 3 of this section, applied to the sequence \(g_1, g_2, \ldots\).)

Step 4: conclusion.

Since \(x^n = (2n+1)^{-1/2} g_n\), linearity of \(T\) gives

\begin{equation*} \sqrt{n} \, \|T(x^n)\|_2 = \sqrt{\frac{n}{2n+1}} \, \|T g_n\|_2 \le \frac{1}{\sqrt{2}} \, \|T g_n\|_2 . \end{equation*}

By Step 3 the right side tends to \(0\), so \(\lim_{n \to \infty} \sqrt{n} \, \|T(x^n)\|_2 = 0\), as desired.

Problem (10C.3)

Suppose \(T\) is a compact operator on a Hilbert space \(V\) and \(f_1, f_2, \ldots\) is a sequence in \(V\) such that \(\lim_{n \to \infty} \langle f_n, g \rangle = 0\) for every \(g \in V\). Prove that \(\lim_{n \to \infty} \|T f_n\| = 0\).

Solution

Step 1: the sequence \(f_1, f_2, \ldots\) is bounded.

For each \(n \in \mathbf{Z}^+\) define \(\varphi_n \colon V \to \mathbf{F}\) by

\begin{equation*} \varphi_n(g) = \langle g, f_n \rangle . \end{equation*}

Then \(\varphi_n\) is a linear functional (the inner product on a Hilbert space is linear in its first slot), and the Cauchy–Schwarz inequality gives \(|\varphi_n(g)| \le \|f_n\| \, \|g\|\), so \(\|\varphi_n\| \le \|f_n\|\). Taking \(g = f_n\) gives \(\varphi_n(f_n) = \|f_n\|^2\), whence \(\|\varphi_n\| \ge \|f_n\|\). Thus \(\|\varphi_n\| = \|f_n\|\), and in particular each \(\varphi_n\) is a bounded linear map from \(V\) to \(\mathbf{F}\).

For each fixed \(g \in V\) we have

\begin{equation*} \varphi_n(g) = \langle g, f_n \rangle = \overline{\langle f_n, g \rangle} \longrightarrow 0 \end{equation*}

as \(n \to \infty\), by hypothesis. A convergent sequence of scalars is bounded, so

\begin{equation*} \sup \{ |\varphi_n(g)| : n \in \mathbf{Z}^+ \} < \infty \qquad \text{for every } g \in V . \end{equation*}

A Hilbert space is a Banach space, so the Principle of Uniform Boundedness (6.86), applied to the family \(\{\varphi_n\}_{n \in \mathbf{Z}^+}\) of bounded linear maps from \(V\) to \(\mathbf{F}\), gives

\begin{equation*} M := \sup_{n \in \mathbf{Z}^+} \|f_n\| = \sup_{n \in \mathbf{Z}^+} \|\varphi_n\| < \infty . \end{equation*}

Thus \(f_1, f_2, \ldots\) is a bounded sequence in \(V\).

Step 2: the compactness argument.

Suppose, seeking a contradiction, that \(\|T f_n\|\) does not converge to \(0\). Then there exist \(\varepsilon > 0\) and a subsequence \(f_{n_1}, f_{n_2}, \ldots\) with

\begin{equation*} \|T f_{n_k}\| \ge \varepsilon \qquad \text{for every } k \in \mathbf{Z}^+ . \end{equation*}

By Step 1 this subsequence is bounded, so compactness of \(T\) (10.66) provides a further subsequence \(f_{n_{k_1}}, f_{n_{k_2}}, \ldots\) and a vector \(u \in V\) with

\begin{equation*} \lim_{j \to \infty} T f_{n_{k_j}} = u . \end{equation*}

The Cauchy–Schwarz inequality implies that the inner product is continuous in its first slot, so

\begin{equation*} \lim_{j \to \infty} \langle T f_{n_{k_j}}, u \rangle = \langle u, u \rangle = \|u\|^2 . \end{equation*}

On the other hand, \(T\) is bounded (10.68), so \(T^*\) exists and \(\langle T f_n, u \rangle = \langle f_n, T^* u \rangle\) for every \(n\). Taking \(g = T^* u\) in the hypothesis gives

\begin{equation*} \lim_{n \to \infty} \langle T f_n, u \rangle = \lim_{n \to \infty} \langle f_n, T^* u \rangle = 0 , \end{equation*}

and the same holds along the subsequence. Comparing the two displayed limits gives \(\|u\|^2 = 0\), hence \(u = 0\).

But then \(\|T f_{n_{k_j}}\| \to \|u\| = 0\), contradicting \(\|T f_{n_{k_j}}\| \ge \varepsilon > 0\) for all \(j\).

This contradiction shows that \(\lim_{n \to \infty} \|T f_n\| = 0\), as desired.

Remark. The hypothesis says precisely that \(f_1, f_2, \ldots\) converges weakly to \(0\); the conclusion is that a compact operator maps weakly convergent sequences to norm convergent ones. Exercise 1 is the special case \(f_n = e_n\), because Bessel’s inequality (8.57) shows that every orthonormal sequence converges weakly to \(0\).

Problem (10C.4)

Suppose \(h \in L^{\infty}(\mathbb{R})\). Define \(M_h \in \mathcal{B}\bigl( L^2(\mathbb{R}) \bigr)\) by \(M_h f = f h\). Prove that if \(\|h\|_{\infty} > 0\), then \(M_h\) is not compact.

Solution

Here \(\lambda\) denotes Lebesgue measure on \(\mathbb{R}\). Note first that \(M_h\) really is a bounded operator on \(L^2(\mathbb{R})\), because \(|f h| \le \|h\|_{\infty} |f|\) almost everywhere and hence \(\|M_h f\|_2 \le \|h\|_{\infty} \|f\|_2\).

Assume \(\|h\|_{\infty} > 0\). We produce a bounded sequence in \(L^2(\mathbb{R})\) whose image under \(M_h\) has no convergent subsequence.

Step 1: a set of positive measure on which \(|h|\) is bounded below.

Let \(c = \tfrac{1}{2} \|h\|_{\infty}\), so \(0 < c < \|h\|_{\infty} < \infty\), and set

\begin{equation*} E = \{ x \in \mathbb{R} : |h(x)| > c \} . \end{equation*}

If \(\lambda(E) = 0\), then \(|h| \le c\) almost everywhere, which would give \(\|h\|_{\infty} \le c < \|h\|_{\infty}\), a contradiction. Hence \(\lambda(E) > 0\).

Since \(E = \bigcup_{n \in \mathbb{Z}} \bigl( E \cap [n, n+1) \bigr)\) is a countable union, countable additivity forces \(\lambda\bigl( E \cap [n_0, n_0 + 1) \bigr) > 0\) for some \(n_0 \in \mathbb{Z}\). Put

\begin{equation*} E_0 = E \cap [n_0, n_0 + 1) , \qquad 0 < \lambda(E_0) \le 1 . \end{equation*}

Step 2: split \(E_0\) into infinitely many disjoint pieces of positive measure.

Define \(\varphi \colon [n_0, n_0+1] \to [0, \infty)\) by

\begin{equation*} \varphi(t) = \lambda\bigl( E_0 \cap (-\infty, t] \bigr) . \end{equation*}

If \(s \le t\) then \(0 \le \varphi(t) - \varphi(s) = \lambda\bigl( E_0 \cap (s, t] \bigr) \le t - s\), so \(\varphi\) is nondecreasing and Lipschitz, hence continuous. Also \(\varphi(n_0) = 0\) and \(\varphi(n_0 + 1) = \lambda(E_0)\).

By the Intermediate Value Theorem, for each \(k \in \mathbf{Z}^+\) there exists \(t_k \in [n_0, n_0+1]\) with

\begin{equation*} \varphi(t_k) = \lambda(E_0) \bigl( 1 - 2^{-k} \bigr) . \end{equation*}

The numbers \(\varphi(t_1) < \varphi(t_2) < \cdots\) are strictly increasing and \(\varphi\) is nondecreasing, so \(t_1 < t_2 < \cdots\). Now set

\begin{equation*} A_1 = E_0 \cap (-\infty, t_1] , \qquad A_k = E_0 \cap (t_{k-1}, t_k] \ \text{ for } k \ge 2 . \end{equation*}

These sets are pairwise disjoint Borel subsets of \(E_0\), and

\begin{equation*} \lambda(A_1) = \tfrac{1}{2} \lambda(E_0) , \qquad \lambda(A_k) = \varphi(t_k) - \varphi(t_{k-1}) = 2^{-k} \lambda(E_0) \ \text{ for } k \ge 2 . \end{equation*}

In particular \(0 < \lambda(A_k) < \infty\) for every \(k\).

Step 3: the bounded sequence.

For \(k \in \mathbf{Z}^+\) define

\begin{equation*} f_k = \frac{\chi_{A_k}}{\sqrt{\lambda(A_k)}} \in L^2(\mathbb{R}) . \end{equation*}

Then \(\|f_k\|_2 = 1\) for every \(k\), so \(f_1, f_2, \ldots\) is a bounded sequence.

Because \(A_k \subseteq E_0 \subseteq E\), we have \(|h| > c\) on \(A_k\), and therefore

\begin{equation*} \|M_h f_k\|_2^2 = \frac{1}{\lambda(A_k)} \int_{A_k} |h|^2 \, d\lambda \ \ge \ \frac{c^2 \lambda(A_k)}{\lambda(A_k)} = c^2 . \end{equation*}

The functions \(M_h f_1, M_h f_2, \ldots\) have pairwise disjoint supports (since \(M_h f_k\) vanishes outside \(A_k\)), so for \(j \ne k\),

\begin{equation*} \|M_h f_j - M_h f_k\|_2^2 = \|M_h f_j\|_2^2 + \|M_h f_k\|_2^2 \ \ge \ 2 c^2 , \end{equation*}

which gives

\begin{equation*} \|M_h f_j - M_h f_k\|_2 \ \ge \ c \sqrt{2} > 0 \qquad \text{whenever } j \ne k . \end{equation*}

Step 4: conclusion.

The inequality above shows that no subsequence of \(M_h f_1, M_h f_2, \ldots\) is a Cauchy sequence, hence no subsequence converges. Since \(f_1, f_2, \ldots\) is bounded, the definition of compactness (10.66) fails for \(M_h\). Therefore \(M_h\) is not compact.

Problem (10C.5)

Suppose \((b_1, b_2, \ldots) \in \ell^{\infty}\). Define \(T \colon \ell^2 \to \ell^2\) by

\begin{equation*} T(a_1, a_2, \ldots) = (a_1 b_1, a_2 b_2, \ldots) . \end{equation*}

Prove that \(T\) is compact if and only if \(\lim_{n \to \infty} b_n = 0\).

Solution

First note that \(T\) does map \(\ell^2\) into \(\ell^2\) and is bounded: if \(a = (a_1, a_2, \ldots) \in \ell^2\), then

\begin{equation*} \|T a\|^2 = \sum_{k=1}^{\infty} |a_k b_k|^2 \le \|b\|_{\infty}^2 \sum_{k=1}^{\infty} |a_k|^2 = \|b\|_{\infty}^2 \|a\|^2 , \end{equation*}

where \(b = (b_1, b_2, \ldots)\). Hence \(T \in \mathcal{B}(\ell^2)\) with \(\|T\| \le \|b\|_{\infty}\). Clearly \(T\) is linear.

Let \(e_1, e_2, \ldots\) denote the standard orthonormal basis of \(\ell^2\): \(e_k\) is the sequence that is \(0\) in every slot except for a \(1\) in the \(k^{\text{th}}\) slot. Note that

\begin{equation*} T e_k = b_k e_k , \qquad \text{so} \qquad \|T e_k\| = |b_k| . \end{equation*}

Suppose \(\lim_{n \to \infty} b_n = 0\); we show \(T\) is compact.

This is Example 10.72; here is the argument. For \(n \in \mathbf{Z}^+\) define \(T_n \colon \ell^2 \to \ell^2\) by

\begin{equation*} T_n(a_1, a_2, \ldots) = (a_1 b_1, \ldots, a_n b_n, 0, 0, \ldots) . \end{equation*}

Each \(T_n\) is a bounded linear operator (by the same estimate as above) whose range is contained in the \(n\)-dimensional subspace \(\operatorname{span}(e_1, \ldots, e_n)\). Hence each \(T_n\) is compact, by 10.67.

For \(a \in \ell^2\) we have \((T - T_n) a = (0, \ldots, 0, a_{n+1} b_{n+1}, a_{n+2} b_{n+2}, \ldots)\), so

\begin{equation*} \|(T - T_n) a\|^2 = \sum_{k > n} |a_k b_k|^2 \le \Bigl( \sup_{k > n} |b_k| \Bigr)^2 \|a\|^2 . \end{equation*}

Thus \(\|T - T_n\| \le \sup_{k > n} |b_k|\). Because \(\lim_{k \to \infty} b_k = 0\), we have \(\sup_{k > n} |b_k| \to 0\) as \(n \to \infty\), and therefore

\begin{equation*} \lim_{n \to \infty} \|T - T_n\| = 0 . \end{equation*}

Since \(\mathcal{C}(\ell^2)\) is a closed subset of \(\mathcal{B}(\ell^2)\) by 10.69(a), and \(T\) is the limit in norm of the compact operators \(T_1, T_2, \ldots\), we conclude that \(T\) is compact.

Suppose \(T\) is compact; we show \(\lim_{n \to \infty} b_n = 0\).

The sequence \(e_1, e_2, \ldots\) is orthonormal in \(\ell^2\). Exercise 1 of this section states that \(\lim_{n \to \infty} T e_n = 0\) for every compact operator \(T\) and every orthonormal sequence. Hence

\begin{equation*} \lim_{n \to \infty} |b_n| = \lim_{n \to \infty} \|T e_n\| = 0 , \end{equation*}

which says exactly that \(\lim_{n \to \infty} b_n = 0\).

(For a direct argument avoiding Exercise 1: if \(b_n \not\to 0\), there exist \(\varepsilon > 0\) and \(n_1 < n_2 < \cdots\) with \(|b_{n_k}| \ge \varepsilon\). The sequence \(e_{n_1}, e_{n_2}, \ldots\) is bounded, but for \(j \ne k\),

\begin{equation*} \|T e_{n_j} - T e_{n_k}\|^2 = \|b_{n_j} e_{n_j} - b_{n_k} e_{n_k}\|^2 = |b_{n_j}|^2 + |b_{n_k}|^2 \ge 2 \varepsilon^2 \end{equation*}

because \(e_{n_j} \perp e_{n_k}\). Hence \(T e_{n_1}, T e_{n_2}, \ldots\) has no Cauchy subsequence and \(T\) is not compact.)

Combining the two implications completes the proof.

Problem (10C.6)

Suppose \(T\) is a bounded operator on a Hilbert space \(V\). Prove that if there exists an orthonormal basis \(\{e_k\}_{k \in \Gamma}\) of \(V\) such that

\begin{equation*} \sum_{k \in \Gamma} \|T e_k\|^2 < \infty , \end{equation*}

then \(T\) is compact.

Solution

Assume \(\{e_k\}_{k \in \Gamma}\) is an orthonormal basis of \(V\) with \(S := \sum_{k \in \Gamma} \|T e_k\|^2 < \infty\). The strategy is to approximate \(T\) in norm by operators with finite-dimensional range and then invoke the closedness of \(\mathcal{C}(V)\).

Step 1: tails of the sum are small.

The terms \(\|T e_k\|^2\) are nonnegative, so (by the first bullet point following 8.53)

\begin{equation*} S = \sup \Bigl\{ \sum_{k \in \Omega} \|T e_k\|^2 : \Omega \text{ is a finite subset of } \Gamma \Bigr\} < \infty . \end{equation*}

Let \(\varepsilon > 0\). Choose a finite set \(\Omega \subseteq \Gamma\) with

\begin{equation*} \sum_{k \in \Omega} \|T e_k\|^2 > S - \varepsilon^2 . \end{equation*}

Then for every finite set \(F \subseteq \Gamma \setminus \Omega\) we have, since \(\Omega \cup F\) is finite and disjointly decomposed,

\begin{equation*} \sum_{k \in F} \|T e_k\|^2 = \sum_{k \in \Omega \cup F} \|T e_k\|^2 - \sum_{k \in \Omega} \|T e_k\|^2 \le S - (S - \varepsilon^2) = \varepsilon^2 . \end{equation*}

Step 2: the approximating operator.

Define \(P_{\Omega} \colon V \to V\) by

\begin{equation*} P_{\Omega} f = \sum_{k \in \Omega} \langle f, e_k \rangle e_k \end{equation*}

(a finite sum). This is linear, and \(\|P_{\Omega} f\|^2 = \sum_{k \in \Omega} |\langle f, e_k \rangle|^2 \le \|f\|^2\) by Bessel’s inequality (8.57), so \(P_{\Omega}\) is bounded. Its range lies in the finite-dimensional subspace spanned by \(\{e_k\}_{k \in \Omega}\). Hence \(T P_{\Omega}\) is a bounded operator with finite-dimensional range, and so \(T P_{\Omega}\) is compact by 10.67.

Step 3: \(\|T - T P_{\Omega}\| \le \varepsilon\).

Let \(f \in V\). By Parseval’s identity 8.63(a),

\begin{equation*} f = \sum_{k \in \Gamma} \langle f, e_k \rangle e_k . \end{equation*}

Splitting off the finite piece indexed by \(\Omega\) gives

\begin{equation*} f - P_{\Omega} f = \sum_{k \in \Gamma \setminus \Omega} \langle f, e_k \rangle e_k , \end{equation*}

where the unordered sum on the right converges. Indeed, if \(g\) denotes the sum over \(\Gamma\) and \(\eta > 0\), pick a finite \(\Omega_0 \subseteq \Gamma\) as in 8.53; then for every finite \(F\) with \(\Omega_0 \setminus \Omega \subseteq F \subseteq \Gamma \setminus \Omega\) the set \(\Omega \cup F\) contains \(\Omega_0\), so

\begin{equation*} \Bigl\| (g - P_{\Omega} f) - \sum_{k \in F} \langle f, e_k \rangle e_k \Bigr\| = \Bigl\| g - \sum_{k \in \Omega \cup F} \langle f, e_k \rangle e_k \Bigr\| < \eta , \end{equation*}

which is exactly the assertion.

Because \(T\) is bounded, hence continuous, applying \(T\) to a convergent unordered sum gives the convergent unordered sum of the images: for finite \(F\),

\begin{equation*} \Bigl\| T(f - P_{\Omega} f) - \sum_{k \in F} \langle f, e_k \rangle T e_k \Bigr\| \le \|T\| \, \Bigl\| (f - P_{\Omega} f) - \sum_{k \in F} \langle f, e_k \rangle e_k \Bigr\| , \end{equation*}

so

\begin{equation*} T(f - P_{\Omega} f) = \sum_{k \in \Gamma \setminus \Omega} \langle f, e_k \rangle \, T e_k . \end{equation*}

Now estimate a finite partial sum. For finite \(F \subseteq \Gamma \setminus \Omega\), the triangle inequality, the Cauchy–Schwarz inequality in \(\mathbb{R}^{F}\), Step 1, and Bessel’s inequality give

\begin{equation*} \begin{aligned} \Bigl\| \sum_{k \in F} \langle f, e_k \rangle T e_k \Bigr\| &\le \sum_{k \in F} |\langle f, e_k \rangle| \, \|T e_k\| \\ &\le \Bigl( \sum_{k \in F} |\langle f, e_k \rangle|^2 \Bigr)^{1/2} \Bigl( \sum_{k \in F} \|T e_k\|^2 \Bigr)^{1/2} \\ &\le \|f\| \cdot \varepsilon . \end{aligned} \end{equation*}

Since \(T(f - P_{\Omega} f)\) is within \(\eta\) of such a partial sum for suitable finite \(F\) and every \(\eta > 0\), we conclude

\begin{equation*} \|(T - T P_{\Omega}) f\| = \|T(f - P_{\Omega} f)\| \le \varepsilon \|f\| . \end{equation*}

As \(f \in V\) was arbitrary, \(\|T - T P_{\Omega}\| \le \varepsilon\).

Step 4: conclusion.

For each \(m \in \mathbf{Z}^+\), apply Steps 1–3 with \(\varepsilon = 1/m\) to obtain a finite set \(\Omega_m \subseteq \Gamma\) with

\begin{equation*} \|T - T P_{\Omega_m}\| \le \tfrac{1}{m} . \end{equation*}

Each \(T P_{\Omega_m}\) is compact, and \(\lim_{m \to \infty} \|T - T P_{\Omega_m}\| = 0\). Since \(\mathcal{C}(V)\) is a closed subset of \(\mathcal{B}(V)\) by 10.69(a), it follows that \(T\) is compact.

Remark. Operators satisfying the hypothesis are called Hilbert–Schmidt operators; Exercise 7 shows that the value of \(\sum_{k \in \Gamma} \|T e_k\|^2\) does not depend on the choice of orthonormal basis.

Problem (10C.7)

Suppose \(T\) is a bounded operator on a Hilbert space \(V\). Prove that if \(\{e_k\}_{k \in \Gamma}\) and \(\{f_j\}_{j \in \Omega}\) are orthonormal bases of \(V\), then

\begin{equation*} \sum_{k \in \Gamma} \|T e_k\|^2 = \sum_{j \in \Omega} \|T f_j\|^2 . \end{equation*}

Solution

All sums below have nonnegative terms, so by the first bullet point following 8.53 each is the supremum of its finite partial sums and takes a value in \([0, \infty]\); the asserted equality is understood in \([0, \infty]\) (in particular one side is infinite exactly when the other is).

Step 1: a Tonelli-type lemma for nonnegative unordered sums.

Claim. If \(a_{k,j} \in [0, \infty)\) for all \(k \in \Gamma\) and \(j \in \Omega\), then

\begin{equation*} \sum_{k \in \Gamma} \; \sum_{j \in \Omega} a_{k,j} \; = \; \sum_{j \in \Omega} \; \sum_{k \in \Gamma} a_{k,j} \end{equation*}

as elements of \([0, \infty]\). Here an outer sum whose terms lie in \([0, \infty]\) is again defined as the supremum of its finite partial sums, with the convention that a sum having a term equal to \(\infty\) equals \(\infty\).

Proof of claim. Write \(L\) for the left side and \(R\) for the right side. Fix a finite set \(G \subseteq \Omega\), and for each \(j \in G\) let \(c_j = \sum_{k \in \Gamma} a_{k,j} \in [0, \infty]\). Choose real numbers \(M_j < c_j\) arbitrarily. Because \(c_j\) is the supremum of the finite partial sums \(\sum_{k \in F^{\prime}} a_{k,j}\) over finite \(F^{\prime} \subseteq \Gamma\) (by the first bullet point following 8.53), for each \(j \in G\) we may pick a finite set \(F_j \subseteq \Gamma\) with \(\sum_{k \in F_j} a_{k,j} > M_j\). Let \(F = \bigcup_{j \in G} F_j\), a finite subset of \(\Gamma\). Since \(F_j \subseteq F\) and all terms are nonnegative, \(\sum_{k \in F} a_{k,j} \ge \sum_{k \in F_j} a_{k,j} > M_j\) for each \(j \in G\); hence

\begin{equation*} \sum_{j \in G} M_j \; \le \; \sum_{j \in G} \sum_{k \in F} a_{k,j} \; = \; \sum_{k \in F} \sum_{j \in G} a_{k,j} \; \le \; \sum_{k \in F} \sum_{j \in \Omega} a_{k,j} \; \le \; L , \end{equation*}

where the middle equality is just a rearrangement of a finite double sum, and the last inequality holds because \(L\) is the supremum of the finite partial sums of the family \(\bigl\{ \sum_{j \in \Omega} a_{k,j} \bigr\}_{k \in \Gamma}\) (and \(L = \infty\) if any of those inner sums is infinite). Because the real numbers \(M_j < c_j\) were arbitrary, taking the supremum over all such choices gives \(\sum_{j \in G} c_j \le L\). Taking the supremum over finite \(G \subseteq \Omega\) gives \(R \le L\). Interchanging the roles of \(\Gamma\) and \(\Omega\) gives \(L \le R\), proving the claim.

Step 2: a formula relating \(T\) and \(T^{*}\).

Let \(\{e_k\}_{k \in \Gamma}\) and \(\{f_j\}_{j \in \Omega}\) be any orthonormal bases of \(V\). Because \(\{f_j\}_{j \in \Omega}\) is an orthonormal basis, Parseval’s identity 8.63(c) applied to the vector \(T e_k\) gives

\begin{equation*} \|T e_k\|^2 = \sum_{j \in \Omega} |\langle T e_k, f_j \rangle|^2 \end{equation*}

for every \(k \in \Gamma\). Summing over \(k\) and applying Step 1 with \(a_{k,j} = |\langle T e_k, f_j \rangle|^2\),

\begin{equation*} \sum_{k \in \Gamma} \|T e_k\|^2 = \sum_{k \in \Gamma} \sum_{j \in \Omega} |\langle T e_k, f_j \rangle|^2 = \sum_{j \in \Omega} \sum_{k \in \Gamma} |\langle T e_k, f_j \rangle|^2 . \end{equation*}

By the definition of the adjoint, \(\langle T e_k, f_j \rangle = \langle e_k, T^{*} f_j \rangle\), so

\begin{equation*} |\langle T e_k, f_j \rangle|^2 = |\langle e_k, T^{*} f_j \rangle|^2 = |\langle T^{*} f_j, e_k \rangle|^2 . \end{equation*}

Because \(\{e_k\}_{k \in \Gamma}\) is an orthonormal basis, Parseval’s identity 8.63(c) applied to \(T^{*} f_j\) gives \(\sum_{k \in \Gamma} |\langle T^{*} f_j, e_k \rangle|^2 = \|T^{*} f_j\|^2\). Therefore

\begin{equation*} \sum_{k \in \Gamma} \|T e_k\|^2 = \sum_{j \in \Omega} \|T^{*} f_j\|^2 \end{equation*}

for every pair of orthonormal bases \(\{e_k\}_{k \in \Gamma}\) and \(\{f_j\}_{j \in \Omega}\) of \(V\). Call this identity \((\star)\).

Step 3: conclusion.

Apply \((\star)\) twice.

Taking both bases in \((\star)\) to be \(\{f_j\}_{j \in \Omega}\) gives

\begin{equation*} \sum_{j \in \Omega} \|T f_j\|^2 = \sum_{j \in \Omega} \|T^{*} f_j\|^2 . \end{equation*}

Taking the first basis to be \(\{e_k\}_{k \in \Gamma}\) and the second to be \(\{f_j\}_{j \in \Omega}\) gives

\begin{equation*} \sum_{k \in \Gamma} \|T e_k\|^2 = \sum_{j \in \Omega} \|T^{*} f_j\|^2 . \end{equation*}

Combining the two displayed equations yields

\begin{equation*} \sum_{k \in \Gamma} \|T e_k\|^2 = \sum_{j \in \Omega} \|T f_j\|^2 , \end{equation*}

as desired.

Remark. The common value is the square of the Hilbert–Schmidt norm of \(T\), and \((\star)\) with \(\{f_j\} = \{e_k\}\) shows in addition that \(\sum_{k \in \Gamma} \|T e_k\|^2 = \sum_{k \in \Gamma} \|T^{*} e_k\|^2\), so \(T\) is a Hilbert–Schmidt operator if and only if \(T^{*}\) is.

Problem (10C.8)

Suppose \(T\) is a bounded operator on a Hilbert space. Prove that \(T\) is compact if and only if \(T^*T\) is compact.

Solution

Let \(V\) be the Hilbert space and let \(T \in \mathcal{B}(V)\).

First suppose \(T\) is compact. Because \(T^*\) is a bounded operator on \(V\), 10.69(b) implies that \(T^*T \in \mathcal{C}(V)\). Thus \(T^*T\) is compact.

Now suppose \(T^*T\) is compact. Let \(f_1, f_2, \ldots\) be a bounded sequence in \(V\), say

\begin{equation*} M = \sup_{n \in \mathbb{Z}^+} \|f_n\| < \infty . \end{equation*}

Because \(T^*T\) is compact, there is an increasing sequence \(n_1 < n_2 < \cdots\) of positive integers such that \(T^*Tf_{n_1}, T^*Tf_{n_2}, \ldots\) converges; in particular that sequence is Cauchy. For all \(j, k \in \mathbb{Z}^+\) we have, writing \(h = f_{n_j} - f_{n_k}\),

\begin{equation*} \|Tf_{n_j} - Tf_{n_k}\|^2 = \langle Th, Th \rangle = \langle T^*Th, h \rangle \le \|T^*Th\| \, \|h\| \le 2M \, \|T^*Tf_{n_j} - T^*Tf_{n_k}\| , \end{equation*}

where the first inequality is the Cauchy-Schwarz inequality and the second uses \(\|h\| \le \|f_{n_j}\| + \|f_{n_k}\| \le 2M\).

Given \(\varepsilon > 0\), choose \(N\) so large that \(\|T^*Tf_{n_j} - T^*Tf_{n_k}\| < \varepsilon^2/(2M+1)\) for all \(j, k \ge N\) (this is possible because \(T^*Tf_{n_1}, T^*Tf_{n_2}, \ldots\) is Cauchy). The displayed inequality then gives \(\|Tf_{n_j} - Tf_{n_k}\| < \varepsilon\) for all \(j, k \ge N\). Hence \(Tf_{n_1}, Tf_{n_2}, \ldots\) is a Cauchy sequence in \(V\), and therefore converges because \(V\) is complete.

Thus every bounded sequence \(f_1, f_2, \ldots\) in \(V\) has a subsequence along which \(T\) produces a convergent sequence, which means \(T\) is compact.

Problem (10C.9)

Prove that if \(T\) is a compact operator on an infinite-dimensional Hilbert space, then \(\|I - T\| \ge 1\).

Solution

Let \(V\) be an infinite-dimensional Hilbert space and let \(T \in \mathcal{C}(V)\).

Because \(V\) is infinite-dimensional, no finite list spans \(V\). Hence we may choose \(f_1, f_2, \ldots\) in \(V\) inductively so that \(f_1 \ne 0\) and \(f_{n+1} \notin \operatorname{span}\{f_1, \ldots, f_n\}\) for each \(n \in \mathbb{Z}^+\); this sequence is linearly independent. Applying the Gram-Schmidt process to it (see the proof of 8.67) produces an orthonormal sequence \(e_1, e_2, \ldots\) in \(V\).

Since \(T\) is compact and \(e_1, e_2, \ldots\) is bounded, there is an increasing sequence \(n_1 < n_2 < \cdots\) of positive integers and there is \(g \in V\) with

\begin{equation*} \lim_{k \to \infty} T e_{n_k} = g . \end{equation*}

We claim \(g = 0\). To see this, let \(h \in V\). Bessel’s inequality (8.57) gives \(\sum_{n \in \mathbb{Z}^+} |\langle h, e_n \rangle|^2 \le \|h\|^2 < \infty\), and hence \(\lim_{n \to \infty} \langle e_n, h \rangle = 0\) (the terms of a convergent series tend to \(0\), and \(\langle e_n, h \rangle = \overline{\langle h, e_n \rangle}\)). Therefore

\begin{equation*} \langle g, h \rangle = \lim_{k \to \infty} \langle T e_{n_k}, h \rangle = \lim_{k \to \infty} \langle e_{n_k}, T^* h \rangle = 0 , \end{equation*}

where the first equality holds because \(Te_{n_k} \to g\) in norm and the inner product is continuous in its first slot. Since \(\langle g, h \rangle = 0\) for every \(h \in V\), taking \(h = g\) gives \(\|g\|^2 = 0\), so \(g = 0\). Thus \(\lim_{k \to \infty} \|Te_{n_k}\| = 0\).

Now for each \(k \in \mathbb{Z}^+\), because \(\|e_{n_k}\| = 1\) we have

\begin{equation*} \|I - T\| \ge \|(I - T)e_{n_k}\| = \|e_{n_k} - Te_{n_k}\| \ge \|e_{n_k}\| - \|Te_{n_k}\| = 1 - \|Te_{n_k}\| . \end{equation*}

Letting \(k \to \infty\) gives \(\|I - T\| \ge 1\), as desired.

Problem (10C.10)

Show that if \(T\) is a surjective but not injective operator on a vector space \(V\), then

\begin{equation*} \operatorname{null} T \subsetneq \operatorname{null} T^2 \subsetneq \operatorname{null} T^3 \subsetneq \cdots . \end{equation*}

Solution

Suppose \(T\) is a surjective but not injective operator on a vector space \(V\). This is the analogue, for surjectivity, of 10.83.

First we verify the inclusions. Suppose \(n \in \mathbb{Z}^+\) and \(f \in \operatorname{null} T^n\). Then

\begin{equation*} T^{n+1} f = T(T^n f) = T 0 = 0 , \end{equation*}

so \(f \in \operatorname{null} T^{n+1}\). Thus \(\operatorname{null} T^n \subseteq \operatorname{null} T^{n+1}\).

Now we show that no inclusion above is an equality. Because \(T\) is not injective, there exists

\begin{equation*} f \in V \quad \text{with} \quad f \ne 0 \quad \text{and} \quad Tf = 0 . \end{equation*}

Fix \(n \in \mathbb{Z}^+\). Because \(T\) is surjective, so is \(T^n\) (a composition of surjective maps is surjective). Hence there exists \(g \in V\) with

\begin{equation*} T^n g = f . \end{equation*}

Then

\begin{equation*} T^{n+1} g = T(T^n g) = Tf = 0 , \end{equation*}

so \(g \in \operatorname{null} T^{n+1}\). However \(T^n g = f \ne 0\), so \(g \notin \operatorname{null} T^n\). Thus \(\operatorname{null} T^n \subsetneq \operatorname{null} T^{n+1}\).

Because \(n \in \mathbb{Z}^+\) was arbitrary, we conclude that

\begin{equation*} \operatorname{null} T \subsetneq \operatorname{null} T^2 \subsetneq \operatorname{null} T^3 \subsetneq \cdots , \end{equation*}

as desired.

As noted in the book’s remark preceding 10.83, this result implies that every surjective operator on a finite-dimensional vector space is injective: a finite-dimensional vector space admits no infinite strictly increasing chain of subspaces, because the dimension increases by at least \(1\) at each step.

Problem (10C.11)

Suppose \(T\) is a compact operator on a Hilbert space and \(\alpha \in \mathbf{F} \setminus \{0\}\).

(a) Prove that \(\operatorname{range}(T - \alpha I)^{m-1} = \operatorname{range}(T - \alpha I)^m\) for some \(m \in \mathbb{Z}^+\).

(b) Prove that \(\operatorname{null}(T - \alpha I)^{n-1} = \operatorname{null}(T - \alpha I)^n\) for some \(n \in \mathbb{Z}^+\).

(c) Show that the smallest positive integer \(m\) that works in (a) equals the smallest positive integer \(n\) that works in (b).

Solution

Throughout, \(V\) denotes the Hilbert space on which \(T\) operates, and we set

\begin{equation*} S = T - \alpha I , \end{equation*}

with the usual convention \(S^0 = I\), so that \(\operatorname{range} S^0 = V\) and \(\operatorname{null} S^0 = \{0\}\).

Preliminary. For each \(k \in \mathbb{Z}^+\) the Binomial Theorem gives

\begin{equation*} S^k = (T - \alpha I)^k = \sum_{j=0}^{k} \binom{k}{j} (-\alpha)^{k-j} T^j = (-\alpha)^k I + C_k , \qquad C_k = \sum_{j=1}^{k} \binom{k}{j} (-\alpha)^{k-j} T^j . \end{equation*}

Every term of \(C_k\) is a scalar times a product of bounded operators one of which is the compact operator \(T\); hence each term is compact by 10.69(b), and \(C_k\) is compact by 10.69(a). Writing \(\beta_k = -(-\alpha)^k\), which is nonzero because \(\alpha \ne 0\), we have

\begin{equation*} S^k = C_k - \beta_k I \quad \text{with } C_k \text{ compact and } \beta_k \ne 0 . \end{equation*}

Consequently, for every \(k \in \mathbb{Z}^+\):

  • \(\operatorname{range} S^k\) is a closed subspace of \(V\), by 10.77;
  • \(\operatorname{null} S^k\) is finite-dimensional, by 10.82.

Both statements also hold trivially for \(k = 0\).

Part (a). For each \(k \in \mathbb{Z}^+\) and each \(f \in V\) we have \(S^k f = S^{k-1}(Sf) \in \operatorname{range} S^{k-1}\), so \(\operatorname{range} S^k \subseteq \operatorname{range} S^{k-1}\). Thus

\begin{equation*} V = \operatorname{range} S^0 \supseteq \operatorname{range} S^1 \supseteq \operatorname{range} S^2 \supseteq \cdots . \end{equation*}

Suppose, for a contradiction, that no \(m \in \mathbb{Z}^+\) works, so that every inclusion above is strict:

\begin{equation*} V \supsetneq \operatorname{range} S \supsetneq \operatorname{range} S^2 \supsetneq \cdots . \end{equation*}

Fix \(n \in \mathbb{Z}^+\). Then \(\operatorname{range} S^n\) and \(\operatorname{range} S^{n+1}\) are closed subspaces of \(V\) (Preliminary), with \(\operatorname{range} S^{n+1}\) a proper subset of \(\operatorname{range} S^n\). Being a closed subset of the complete metric space \(V\), the space \(U = \operatorname{range} S^n\) is itself a Hilbert space [by 6.16(b)], and \(\operatorname{range} S^{n+1}\) is a closed subspace of \(U\). Pick \(f \in U \setminus \operatorname{range} S^{n+1}\) and use 8.43 inside \(U\) to write \(f = g + h\) with \(g \in \operatorname{range} S^{n+1}\) and \(h \in U\) orthogonal to \(\operatorname{range} S^{n+1}\). Then \(h \ne 0\) (otherwise \(f = g \in \operatorname{range} S^{n+1}\)), so setting \(f_n = h / \|h\|\) gives

\begin{equation*} f_n \in \operatorname{range} S^n \cap \bigl(\operatorname{range} S^{n+1}\bigr)^{\perp} \quad \text{with } \|f_n\| = 1 , \end{equation*}

where the orthogonal complement may equally well be taken in \(U\) or in \(V\), since \(\operatorname{range} S^{n+1} \subseteq U\). Now suppose \(j, k \in \mathbb{Z}^+\) with \(j < k\). Because \(T = S + \alpha I\),

\begin{equation*} T f_j - T f_k = \bigl( S f_j - S f_k - \alpha f_k \bigr) + \alpha f_j . \end{equation*}

Since \(f_j \in \operatorname{range} S^j\) we have \(Sf_j \in \operatorname{range} S^{j+1}\); since \(f_k \in \operatorname{range} S^k \subseteq \operatorname{range} S^{j+1}\) (using \(k \ge j+1\)) we have \(f_k \in \operatorname{range} S^{j+1}\) and \(Sf_k \in \operatorname{range} S^{k+1} \subseteq \operatorname{range} S^{j+1}\). Thus the parenthesized sum, call it \(u\), lies in \(\operatorname{range} S^{j+1}\), and \(f_j \perp \operatorname{range} S^{j+1}\) by choice of \(f_j\); hence \(u \perp \alpha f_j\). The Pythagorean identity now gives

\begin{equation*} \|Tf_j - Tf_k\|^2 = \|u\|^2 + |\alpha|^2 \|f_j\|^2 \ge |\alpha|^2 , \qquad \text{so} \qquad \|Tf_j - Tf_k\| \ge |\alpha| > 0 . \end{equation*}

Hence no subsequence of \(Tf_1, Tf_2, \ldots\) is Cauchy, so \(Tf_1, Tf_2, \ldots\) has no convergent subsequence, even though \(f_1, f_2, \ldots\) is bounded. This contradicts the compactness of \(T\) and proves (a).

Part (b). From \(S^{k-1}f = 0 \Rightarrow S^k f = S(S^{k-1}f) = 0\) we get

\begin{equation*} \{0\} = \operatorname{null} S^0 \subseteq \operatorname{null} S^1 \subseteq \operatorname{null} S^2 \subseteq \cdots , \end{equation*}

and each \(\operatorname{null} S^k\) is closed because \(S^k\) is continuous (a bounded operator, being a polynomial in \(T\)). Suppose, for a contradiction, that no \(n \in \mathbb{Z}^+\) works, so every inclusion above is strict. Fix \(n \in \mathbb{Z}^+\). Then \(\operatorname{null} S^{n}\) is a closed subspace of \(V\), hence a Hilbert space [by 6.16(b)], and \(\operatorname{null} S^{n-1}\) is a proper closed subspace of it. Exactly as in part (a) - pick \(f \in \operatorname{null} S^n \setminus \operatorname{null} S^{n-1}\), split it by 8.43 inside \(\operatorname{null} S^n\), and normalize the component orthogonal to \(\operatorname{null} S^{n-1}\) - we obtain

\begin{equation*} e_n \in \operatorname{null} S^n \cap \bigl(\operatorname{null} S^{n-1}\bigr)^{\perp} \quad \text{with } \|e_n\| = 1 . \end{equation*}

Suppose \(j, k \in \mathbb{Z}^+\) with \(j < k\). Then

\begin{equation*} T e_k - T e_j = \bigl( S e_k - S e_j - \alpha e_j \bigr) + \alpha e_k . \end{equation*}

Now \(Se_k \in \operatorname{null} S^{k-1}\) because \(S^{k-1}(Se_k) = S^k e_k = 0\); also \(e_j \in \operatorname{null} S^{j} \subseteq \operatorname{null} S^{k-1}\) since \(j \le k-1\), and likewise \(Se_j \in \operatorname{null} S^{j-1} \subseteq \operatorname{null} S^{k-1}\). Thus the parenthesized sum, call it \(v\), lies in \(\operatorname{null} S^{k-1}\), and \(e_k \perp \operatorname{null} S^{k-1}\) by choice of \(e_k\); hence \(v \perp \alpha e_k\). The Pythagorean identity now gives

\begin{equation*} \|Te_k - Te_j\|^2 = \|v\|^2 + |\alpha|^2 \|e_k\|^2 \ge |\alpha|^2 , \qquad \text{so} \qquad \|Te_k - Te_j\| \ge |\alpha| > 0 , \end{equation*}

so no subsequence of \(Te_1, Te_2, \ldots\) is Cauchy and hence \(Te_1, Te_2, \ldots\) has no convergent subsequence, although \(e_1, e_2, \ldots\) is bounded. This contradicts the compactness of \(T\) and proves (b).

Part (c). For \(k \in \{0, 1, 2, \ldots\}\) put

\begin{equation*} d_k = \dim \operatorname{null} S^k , \qquad c_k = \dim \bigl( \operatorname{range} S^k \bigr)^{\perp} . \end{equation*}

Step 1: \(d_k = c_k < \infty\) for every \(k\). For \(k = 0\) both are \(0\). For \(k \in \mathbb{Z}^+\), use \(S^k = C_k - \beta_k I\) from the Preliminary. By 10.13(a),

\begin{equation*} \bigl( \operatorname{range} S^k \bigr)^{\perp} = \operatorname{null} (S^k)^* = \operatorname{null} \bigl( C_k^{\,*} - \overline{\beta_k} I \bigr) , \end{equation*}

and \(C_k^{\,*}\) is compact by 10.73, so this space is finite-dimensional by 10.82; also \(\operatorname{null} S^k = \operatorname{null}(C_k - \beta_k I)\) is finite-dimensional by 10.82. Finally, 10.91 applied to the compact operator \(C_k\) and the nonzero scalar \(\beta_k\) gives

\begin{equation*} d_k = \dim \operatorname{null}(C_k - \beta_k I) = \dim \operatorname{null}\bigl(C_k^{\,*} - \overline{\beta_k} I\bigr) = c_k . \end{equation*}

Step 2: for each \(k \in \mathbb{Z}^+\),

\begin{equation*} \operatorname{range} S^{k-1} = \operatorname{range} S^{k} \iff c_{k-1} = c_k . \end{equation*}

Indeed, \(\operatorname{range} S^k \subseteq \operatorname{range} S^{k-1}\) implies \((\operatorname{range} S^{k-1})^{\perp} \subseteq (\operatorname{range} S^{k})^{\perp}\), and these are finite-dimensional by Step 1; hence equality of their dimensions forces \((\operatorname{range} S^{k-1})^{\perp} = (\operatorname{range} S^{k})^{\perp}\). Taking orthogonal complements and using that both ranges are closed (Preliminary), together with 8.41, gives \(\operatorname{range} S^{k-1} = \operatorname{range} S^{k}\). The reverse implication is obvious. Likewise, since \(\operatorname{null} S^{k-1} \subseteq \operatorname{null} S^{k}\) with both finite-dimensional,

\begin{equation*} \operatorname{null} S^{k-1} = \operatorname{null} S^{k} \iff d_{k-1} = d_k . \end{equation*}

Step 3: conclusion. Let \(m\) be the smallest positive integer with \(\operatorname{range} S^{m-1} = \operatorname{range} S^m\) and let \(n\) be the smallest positive integer with \(\operatorname{null} S^{n-1} = \operatorname{null} S^n\); these exist by (a) and (b). By Step 2, \(m\) is the smallest positive integer \(k\) with \(c_{k-1} = c_k\), and \(n\) is the smallest positive integer \(k\) with \(d_{k-1} = d_k\). Because \(d_k = c_k\) for every \(k\) by Step 1, these two smallest integers coincide. Hence \(m = n\).

Remark. The number \(m = n\) is the length of the chain: for \(k < m\) the inclusions \(\operatorname{range} S^{k} \supsetneq \operatorname{range} S^{k+1}\) and \(\operatorname{null} S^{k} \subsetneq \operatorname{null} S^{k+1}\) are strict, while from \(m\) onward both chains are constant. Constancy from \(m\) onward follows because \(\operatorname{range} S^{m-1} = \operatorname{range} S^m\) implies \(\operatorname{range} S^{m} = S(\operatorname{range} S^{m-1}) = S(\operatorname{range} S^{m}) = \operatorname{range} S^{m+1}\), and because \(\operatorname{null} S^{n-1} = \operatorname{null} S^n\) implies that \(S^{n+1}f = 0\) forces \(Sf \in \operatorname{null} S^n = \operatorname{null} S^{n-1}\), that is, \(S^n f = 0\).

Problem (10C.12)

Prove that if \(f : [0, 1] \to \mathbf{F}\) is a continuous function, then there exists a continuous function \(g : [0, 1] \to \mathbf{F}\) such that

\begin{equation*} f(x) = g(x) + \int_0^x g \end{equation*}

for all \(x \in [0, 1]\).

Solution

Let \(\mathcal{V}\) denote the Volterra operator on \(L^2([0,1])\), defined by \((\mathcal{V}h)(x) = \int_0^x h\) (see Example 10.15). Recall that \(\mathcal{V}\) is compact (see the comment after the proof of 10.70) and that \(\operatorname{sp}(\mathcal{V}) = \{0\}\) (Example 10.89).

Existence of \(g\) in \(L^2\). Since \(-1 \ne 0\) and \(\operatorname{sp}(\mathcal{V}) = \{0\}\), we have \(-1 \notin \operatorname{sp}(\mathcal{V})\), which by definition of the spectrum means that

\begin{equation*} \mathcal{V} - (-1)I = I + \mathcal{V} \end{equation*}

is an invertible element of \(\mathcal{B}\bigl(L^2([0,1])\bigr)\). (Alternatively: \(-1\) is not an eigenvalue of \(\mathcal{V}\), as the computation in Example 10.89 shows, and the Fredholm Alternative 10.85 then gives \(-1 \notin \operatorname{sp}(\mathcal{V})\).)

The continuous function \(f\) lies in \(L^2([0,1])\), because \(f\) is bounded on the compact interval \([0,1]\). Hence there exists \(g_0 \in L^2([0,1])\) with

\begin{equation*} g_0 + \mathcal{V} g_0 = f \qquad \text{in } L^2([0,1]) . \end{equation*}

Upgrading to a continuous solution. Because \([0,1]\) has measure \(1\), the Cauchy-Schwarz inequality gives \(\int_{[0,1]} |g_0| = \int_{[0,1]} |g_0| \cdot 1 \le \|g_0\|_2 \, \|1\|_2 = \|g_0\|_2 < \infty\), so \(L^2([0,1]) \subseteq L^1([0,1])\) and in particular \(g_0\) is integrable. Define

\begin{equation*} g(x) = f(x) - \int_0^x g_0 \qquad \text{for } x \in [0,1] . \end{equation*}

The function \(x \mapsto \int_0^x g_0 = \int_{[0,1]} \chi_{[0,x]} g_0\) is continuous on \([0,1]\): if \(x_n \to x\) in \([0,1]\), then

\begin{equation*} \Bigl| \int_0^{x_n} g_0 - \int_0^{x} g_0 \Bigr| \le \int_{[0,1]} \bigl| \chi_{[0,x_n]} - \chi_{[0,x]} \bigr| \, |g_0| \longrightarrow 0 \end{equation*}

by the Dominated Convergence Theorem (3.31), the integrand being dominated by \(|g_0| \in L^1\) and tending to \(0\) pointwise except possibly at \(x\) itself. Since \(f\) is continuous, \(g\) is continuous on \([0,1]\).

Finally, the equation \(g_0 + \mathcal{V}g_0 = f\) says exactly that \(g_0(x) = f(x) - \int_0^x g_0\) for almost every \(x \in [0,1]\); that is, \(g_0 = g\) almost everywhere. Consequently \(\int_0^x g_0 = \int_0^x g\) for every \(x \in [0,1]\), and therefore for every \(x \in [0,1]\),

\begin{equation*} g(x) + \int_0^x g = g(x) + \int_0^x g_0 = f(x) , \end{equation*}

by the definition of \(g\). This is the desired identity, valid at every point of \([0,1]\) (not merely almost everywhere).

An explicit formula. It is worth recording that one can write \(g\) down and verify the conclusion by elementary calculus, with no operator theory. Put

\begin{equation*} h(x) = \int_0^x e^{t-x} f(t) \, dt = e^{-x} \int_0^x e^{t} f(t) \, dt , \qquad g(x) = f(x) - h(x) . \end{equation*}

Since \(f\) is continuous, the Fundamental Theorem of Calculus shows that \(h\) is differentiable with

\begin{equation*} h^{\prime}(x) = -e^{-x} \int_0^x e^t f(t) \, dt + e^{-x} e^x f(x) = f(x) - h(x) , \end{equation*}

and \(h(0) = 0\). Hence \(\int_0^x (f - h) = \int_0^x h^{\prime} = h(x) - h(0) = h(x)\), so

\begin{equation*} g(x) + \int_0^x g = f(x) - h(x) + \int_0^x (f - h) = f(x) - h(x) + h(x) = f(x) \end{equation*}

for all \(x \in [0,1]\). The function \(g = f - h\) is continuous, completing this second proof.

Problem (10C.13)

Suppose \(S\) is a bounded invertible operator on a Hilbert space \(V\) and \(T\) is a compact operator on \(V\).

(a) Prove that \(S + T\) has closed range.

(b) Prove that \(S + T\) is injective if and only if \(S + T\) is surjective.

(c) Prove that \(\operatorname{null}(S + T)\) and \(\operatorname{null}(S^* + T^*)\) are finite-dimensional.

(d) Prove that \(\dim \operatorname{null}(S + T) = \dim \operatorname{null}(S^* + T^*)\).

(e) Prove that there exists \(R \in \mathcal{B}(V)\) such that \(\operatorname{range} R\) is finite-dimensional and \(S + T + R\) is invertible.

Solution

Set up. Because \(S\) is invertible in \(\mathcal{B}(V)\), the operators

\begin{equation*} K = S^{-1}T \qquad \text{and} \qquad L = T S^{-1} \end{equation*}

are compact by 10.69(b). Note the two factorizations

\begin{equation*} S + T = S(I + K) = (I + L)S . \end{equation*}

Also \(S^*\) is invertible, with \((S^*)^{-1} = (S^{-1})^*\), and \(T^*\) is compact by 10.73.

Here \(S^{-1}\) is bounded: by 10.18 an invertible operator is a linear bijection of \(V\) onto \(V\), and the Bounded Inverse Theorem (6.83) makes the inverse of a bounded invertible operator on a Banach space bounded, as the book notes after 10.18. So \(S\) is a homeomorphism of \(V\) onto \(V\), a fact we use repeatedly. Note also that \(K\) and \(L\) are indeed bounded, being products of bounded operators.

Part (a). By 10.77 applied to the compact operator \(K\) and the scalar \(\alpha = -1 \ne 0\), the subspace

\begin{equation*} \operatorname{range}(I + K) = \operatorname{range}\bigl(K - (-1)I\bigr) \end{equation*}

is closed in \(V\). Now

\begin{equation*} \operatorname{range}(S + T) = \operatorname{range}\bigl(S(I+K)\bigr) = S\bigl(\operatorname{range}(I+K)\bigr) , \end{equation*}

the second equality because \(S\) is applied to all of \(\operatorname{range}(I+K)\). Finally, if \(U \subseteq V\) is closed then \(S(U) = (S^{-1})^{-1}(U)\) is closed, being the inverse image of a closed set under the continuous map \(S^{-1}\). Hence \(\operatorname{range}(S+T)\) is closed.

Part (b). Because \(S\) is injective, \(S + T = S(I+K)\) is injective if and only if \(I + K\) is injective; because \(S\) is surjective, \(S + T = S(I+K)\) is surjective if and only if \(I + K\) is surjective. So it suffices to show that \(I + K\) is injective if and only if \(I + K\) is surjective.

Apply the Fredholm Alternative (10.85) to the compact operator \(K\) with \(\alpha = -1\). It says that the three conditions

\begin{equation*} -1 \in \operatorname{sp}(K), \qquad -1 \text{ is an eigenvalue of } K, \qquad K - (-1)I \text{ is not surjective} \end{equation*}

are equivalent. The middle condition says precisely that \(I + K\) is not injective, and the third says precisely that \(I + K\) is not surjective. Hence \(I + K\) fails to be injective exactly when it fails to be surjective, that is, \(I + K\) is injective if and only if \(I + K\) is surjective.

Part (c). Since \(S\) is injective,

\begin{equation*} \operatorname{null}(S + T) = \operatorname{null}\bigl(S(I+K)\bigr) = \operatorname{null}(I + K) = \operatorname{null}\bigl(K - (-1)I\bigr) , \end{equation*}

which is finite-dimensional by 10.82 (applied to the compact operator \(K\) and \(\alpha = -1 \ne 0\)).

For the second space, note that \(S^* + T^* = (S+T)^*\), that \(S^*\) is invertible, and that \(T^*\) is compact by 10.73. So the argument of the previous paragraph, applied with \(S\) replaced by \(S^*\) and \(T\) replaced by \(T^*\), shows that

\begin{equation*} \operatorname{null}(S^* + T^*) = \operatorname{null}\bigl(I + (S^*)^{-1}T^*\bigr) \end{equation*}

is finite-dimensional, since \((S^*)^{-1}T^*\) is compact by 10.69(b).

Part (d). Use the other factorization \(S + T = (I + L)S\) with \(L = TS^{-1}\) compact.

First, \(\dim \operatorname{null}(S+T) = \dim \operatorname{null}(I + L)\). Indeed, for \(f \in V\) we have \(f \in \operatorname{null}(I+L)\) if and only if \((I+L)f = 0\), and since \(S\) is a bijection we may write \(f = Sh\) with \(h = S^{-1}f\); then \((I+L)Sh = (S+T)h\), so

\begin{equation*} \operatorname{null}(I+L) = S\bigl( \operatorname{null}(S+T) \bigr) . \end{equation*}

As \(S\) is a linear isomorphism of \(V\) onto \(V\), it carries the subspace \(\operatorname{null}(S+T)\) isomorphically onto \(\operatorname{null}(I+L)\), so these two subspaces have the same dimension.

Second, \(\dim \operatorname{null}(I+L) = \dim \operatorname{null}(I + L^*)\). This is 10.91 applied to the compact operator \(L\) with \(\alpha = -1\): since \(\overline{-1} = -1\),

\begin{equation*} \dim \operatorname{null}\bigl(L - (-1)I\bigr) = \dim \operatorname{null}\bigl(L^* - (-1)I\bigr) . \end{equation*}

Third, \(\operatorname{null}(I + L^*) = \operatorname{null}(S^* + T^*)\). Indeed \((S+T)^* = \bigl((I+L)S\bigr)^* = S^*(I+L)^* = S^*(I + L^*)\), and \(S^*\) is injective, so

\begin{equation*} \operatorname{null}(S^* + T^*) = \operatorname{null}\bigl(S^*(I+L^*)\bigr) = \operatorname{null}(I + L^*) . \end{equation*}

Combining the three displays gives \(\dim \operatorname{null}(S+T) = \dim \operatorname{null}(S^*+T^*)\), as desired.

Part (e). Let \(N = \operatorname{null}(S+T)\) and let \(d = \dim N\), which is finite by (c). By (a), \(\operatorname{range}(S+T)\) is closed, and by 10.13(a),

\begin{equation*} \bigl(\operatorname{range}(S+T)\bigr)^{\perp} = \operatorname{null}\bigl((S+T)^*\bigr) = \operatorname{null}(S^* + T^*) , \end{equation*}

which by (d) has dimension \(d\) as well.

If \(d = 0\), take \(R = 0\): then \(S+T\) is injective, hence also surjective by (b), hence a bijection of \(V\) onto \(V\) and thus invertible (10.18), with bounded inverse by 6.83; and \(\operatorname{range} R = \{0\}\) is finite-dimensional.

Now suppose \(d \ge 1\). Apply the Gram-Schmidt process to a basis of \(N\) to get an orthonormal basis \(e_1, \ldots, e_d\) of \(N\), and likewise let \(u_1, \ldots, u_d\) be an orthonormal basis of \(\bigl(\operatorname{range}(S+T)\bigr)^{\perp}\). Define \(R : V \to V\) by

\begin{equation*} Rf = \sum_{j=1}^{d} \langle f, e_j \rangle u_j . \end{equation*}

Then \(R\) is linear, and because \(u_1, \ldots, u_d\) is orthonormal the Pythagorean identity and Bessel’s inequality (8.57) give \(\|Rf\|^2 = \sum_{j=1}^d |\langle f, e_j\rangle|^2 \le \|f\|^2\), so \(R \in \mathcal{B}(V)\) with \(\|R\| \le 1\); moreover \(\operatorname{range} R \subseteq \operatorname{span}\{u_1, \ldots, u_d\}\) is finite-dimensional.

\(S+T+R\) is injective: suppose \(f \in V\) and \((S+T)f + Rf = 0\). Now \((S+T)f \in \operatorname{range}(S+T)\) while \(Rf \in \bigl(\operatorname{range}(S+T)\bigr)^{\perp}\), so these two vectors are orthogonal; their sum being \(0\) forces (by the Pythagorean identity) \((S+T)f = 0\) and \(Rf = 0\). From \(Rf = 0\) and the linear independence of \(u_1, \ldots, u_d\) we get \(\langle f, e_j \rangle = 0\) for \(j = 1, \ldots, d\). From \((S+T)f = 0\) we get \(f \in N\), whence \(f = \sum_{j=1}^d \langle f, e_j \rangle e_j = 0\).

\(S+T+R\) is surjective: \(R\) is a bounded operator with finite-dimensional range, so \(R\) is compact by 10.67, and therefore \(T + R\) is compact by 10.69(a). Applying (b) with \(T\) replaced by the compact operator \(T + R\), we conclude that \(S + (T+R)\) is surjective, because it is injective.

Thus \(S+T+R\) is a bijection of \(V\) onto \(V\), hence invertible by 10.18; it is bounded, so its inverse is bounded too by 6.83. Since \(\operatorname{range} R\) is finite-dimensional, the proof is complete.

Problem (10C.14)

Suppose \(T\) is a compact operator on a Hilbert space \(V\). Prove that \(\operatorname{range} T\) is a separable subspace of \(V\).

Solution

For \(n \in \mathbb{Z}^+\) let \(B_n = \{ f \in V : \|f\| \le n \}\). Because every \(f \in V\) lies in \(B_n\) for all \(n \ge \|f\|\),

\begin{equation*} \operatorname{range} T = \bigcup_{n=1}^{\infty} T(B_n) . \end{equation*}

(That \(\operatorname{range} T\) is a subspace of \(V\) is immediate from linearity of \(T\); the point of the exercise is separability, i.e. the existence of a countable dense subset of \(\operatorname{range} T\) in the metric inherited from \(V\).)

Step 1: each \(T(B_n)\) is totally bounded. Fix \(n \in \mathbb{Z}^+\) and suppose \(T(B_n)\) is not totally bounded. Then there exists \(\varepsilon > 0\) such that no finite collection of open balls of radius \(\varepsilon\) centered at points of \(T(B_n)\) covers \(T(B_n)\). Construct inductively a sequence \(f_1, f_2, \ldots\) in \(B_n\) as follows: set \(f_1 = 0\), which lies in \(B_n\); having chosen \(f_1, \ldots, f_k\), the balls of radius \(\varepsilon\) about \(Tf_1, \ldots, Tf_k\) do not cover \(T(B_n)\), so there exists \(f_{k+1} \in B_n\) with

\begin{equation*} \|Tf_{k+1} - Tf_i\| \ge \varepsilon \quad \text{for } i = 1, \ldots, k . \end{equation*}

The resulting sequence satisfies \(\|Tf_j - Tf_k\| \ge \varepsilon\) whenever \(j \ne k\), so \(Tf_1, Tf_2, \ldots\) has no Cauchy subsequence and hence no convergent subsequence. But \(f_1, f_2, \ldots\) is a bounded sequence in \(V\), so this contradicts the compactness of \(T\) (10.66). Therefore \(T(B_n)\) is totally bounded.

Step 2: a totally bounded subset of a metric space is separable. Let \(A\) be totally bounded. For each \(k \in \mathbb{Z}^+\) choose a finite set \(F_{k} \subseteq A\) such that every point of \(A\) is within distance \(1/k\) of some point of \(F_{k}\). Then \(D = \bigcup_{k=1}^{\infty} F_{k}\) is a countable subset of \(A\), and \(D\) is dense in \(A\): given \(a \in A\) and \(\varepsilon > 0\), pick \(k\) with \(1/k < \varepsilon\) and then a point of \(F_k\) within \(1/k < \varepsilon\) of \(a\).

Applying Steps 1 and 2, for each \(n \in \mathbb{Z}^+\) there is a countable set \(D_n \subseteq T(B_n)\) that is dense in \(T(B_n)\).

Step 3: conclusion. Let \(D = \bigcup_{n=1}^{\infty} D_n\), a countable union of countable sets, hence countable, and \(D \subseteq \operatorname{range} T\). We claim \(D\) is dense in \(\operatorname{range} T\). Indeed, suppose \(g \in \operatorname{range} T\) and \(\varepsilon > 0\). Write \(g = Tf\) for some \(f \in V\), and choose \(n \in \mathbb{Z}^+\) with \(n \ge \|f\|\), so that \(g \in T(B_n)\). Since \(D_n\) is dense in \(T(B_n)\), there exists \(h \in D_n \subseteq D\) with \(\|g - h\| < \varepsilon\).

Thus \(D\) is a countable dense subset of \(\operatorname{range} T\), so \(\operatorname{range} T\) is a separable subspace of \(V\).

Problem (10C.15)

Suppose \(T\) is a compact operator on a Hilbert space \(V\) and \(e_1, e_2, \ldots\) is an orthonormal basis of \(\operatorname{range} T\). Let \(P_n\) denote the orthogonal projection of \(V\) onto \(\operatorname{span}\{e_1, \ldots, e_n\}\).

(a) Prove that \(\lim_{n \to \infty} \|T - P_n T\| = 0\).

(b) Prove that a bounded operator on a Hilbert space \(V\) is compact if and only if it is the limit in \(\mathcal{B}(V)\) of a sequence of bounded operators with finite-dimensional range.

Solution

Throughout, let \(U = \overline{\operatorname{range} T}\); thus \(U\) is a closed subspace of \(V\) and hence a Hilbert space. Saying that the orthonormal sequence \(e_1, e_2, \ldots\) is an orthonormal basis of \(\operatorname{range} T\) means that \(\overline{\operatorname{span}\{e_1, e_2, \ldots\}} = U\), in other words that \(e_1, e_2, \ldots\) is an orthonormal basis of the Hilbert space \(U\).

(a) First note that \(\|I - P_n\| \le 1\) for each \(n\): if \(f \in V\), then \(f - P_n f\) is orthogonal to \(P_n f\), so the Pythagorean identity gives

\begin{equation*} \|f\|^2 = \|P_n f\|^2 + \|f - P_n f\|^2 \ge \|f - P_n f\|^2 . \end{equation*}

Next note that

\begin{equation*} \lim_{n \to \infty} \|h - P_n h\| = 0 \qquad \text{for every } h \in U . \end{equation*}

Indeed, because \(e_1, e_2, \ldots\) is an orthonormal basis of \(U\) and \(h \in U\), we have \(h = \sum_{k=1}^{\infty} \langle h, e_k \rangle e_k\), while \(P_n h = \sum_{k=1}^{n} \langle h, e_k \rangle e_k\) by 8.71; hence \(\|h - P_n h\|^2 = \sum_{k=n+1}^{\infty} |\langle h, e_k \rangle|^2\), which tends to \(0\) as \(n \to \infty\) because the series \(\sum_{k} |\langle h, e_k \rangle|^2\) converges (it equals \(\|h\|^2\) by Parseval’s identity).

Now suppose, seeking a contradiction, that \(\|T - P_n T\|\) does not converge to \(0\). Then there exist \(\varepsilon > 0\) and integers \(n_1 < n_2 < \cdots\) such that \(\|T - P_{n_j} T\| > \varepsilon\) for every \(j\), and hence for each \(j\) there exists \(f_j \in V\) with \(\|f_j\| \le 1\) and

\begin{equation*} \|(I - P_{n_j}) T f_j\| > \varepsilon . \end{equation*}

The sequence \(f_1, f_2, \ldots\) is bounded, so the compactness of \(T\) provides a subsequence \(f_{j_1}, f_{j_2}, \ldots\) such that \(T f_{j_1}, T f_{j_2}, \ldots\) converges to some \(g \in V\). Because each \(T f_{j_i}\) lies in \(\operatorname{range} T\) and \(U\) is closed, we have \(g \in U\). For each \(i\),

\begin{equation*} \varepsilon < \|(I - P_{n_{j_i}}) T f_{j_i}\| \le \|(I - P_{n_{j_i}})(T f_{j_i} - g)\| + \|(I - P_{n_{j_i}}) g\| \le \|T f_{j_i} - g\| + \|g - P_{n_{j_i}} g\| , \end{equation*}

where the last inequality uses \(\|I - P_n\| \le 1\). As \(i \to \infty\) the first term tends to \(0\) by the choice of the subsequence and the second term tends to \(0\) because \(g \in U\) (note that \(n_{j_i} \to \infty\)). This contradicts the strict inequality above. Hence \(\lim_{n \to \infty} \|T - P_n T\| = 0\).

(b) First suppose \(S_1, S_2, \ldots\) are bounded operators on \(V\), each with finite-dimensional range, and \(\lim_{n \to \infty} \|S - S_n\| = 0\) for some \(S \in \mathcal{B}(V)\). Each \(S_n\) is compact by 10.67, and \(\mathcal{C}(V)\) is a closed subset of \(\mathcal{B}(V)\) by 10.69(a). Hence \(S\) is compact.

Conversely, suppose \(T \in \mathcal{B}(V)\) is compact. We first check that \(U = \overline{\operatorname{range} T}\) is separable. Let \(B = \{f \in V : \|f\| \le 1\}\). Compactness of \(T\) means that every sequence in \(T(B)\) has a subsequence converging in \(V\). The same then holds for sequences in \(\overline{T(B)}\): if \(x_1, x_2, \ldots \in \overline{T(B)}\), choose \(y_n \in T(B)\) with \(\|x_n - y_n\| < 1/n\), take a convergent subsequence of \(y_1, y_2, \ldots\), and note that the corresponding subsequence of \(x_1, x_2, \ldots\) converges to the same limit, which lies in the closed set \(\overline{T(B)}\). Thus \(\overline{T(B)}\) is a sequentially compact metric space, hence compact, hence separable (a compact metric space is covered, for each \(m \in \mathbf{Z}^+\), by finitely many balls of radius \(1/m\), and the centers of all these balls form a countable dense set). For \(m \in \mathbf{Z}^+\) we have \(T(mB) = m \, T(B)\), which is thus also separable, and

\begin{equation*} \operatorname{range} T = \bigcup_{m=1}^{\infty} T(mB) \end{equation*}

is a countable union of separable sets and hence is separable; therefore its closure \(U\) is separable.

If \(U\) is finite-dimensional, then \(T\) itself has finite-dimensional range and the constant sequence \(T, T, \ldots\) exhibits \(T\) as such a limit. So assume \(U\) is infinite-dimensional. Because \(U\) is a separable Hilbert space, 8.67 provides an orthonormal basis \(\{e_k\}_{k \in \Gamma}\) of \(U\). This family is countable: distinct elements \(e\), \(e^{\prime}\) of an orthonormal family satisfy \(\|e - e^{\prime}\| = \sqrt{2}\) by the Pythagorean identity, so the open balls of radius \(\frac{1}{2}\) centered at the \(e_k\) are pairwise disjoint, and each must contain a point of a fixed countable dense subset of \(U\); distinct indices therefore give distinct points of that countable set. Because \(U\) is infinite-dimensional, \(\Gamma\) is infinite, so we may list the family as a sequence \(e_1, e_2, \ldots\). Let \(P_n\) be the orthogonal projection of \(V\) onto \(\operatorname{span}\{e_1, \ldots, e_n\}\) (a finite-dimensional, hence closed, subspace). The proof of (a) used only that \(e_1, e_2, \ldots\) is an orthonormal basis of \(U\), so it applies here and gives \(\lim_{n \to \infty} \|T - P_n T\| = 0\). Each \(P_n T\) is bounded with \(\operatorname{range}(P_n T) \subseteq \operatorname{span}\{e_1, \ldots, e_n\}\), which is finite-dimensional. Thus \(T\) is the limit in \(\mathcal{B}(V)\) of a sequence of bounded operators with finite-dimensional range.

Problem (10C.16)

Prove that if \(T\) is a compact operator on a Hilbert space \(V\), then there exists a sequence \(S_1, S_2, \ldots\) of invertible operators on \(V\) such that \(\lim_{n \to \infty} \|T - S_n\| = 0\).

Solution

The key point is that the spectrum of a compact operator is a countable set, while \(\mathbf{F}\) is not countable; hence we can shrink \(T\) by tiny multiples of the identity that avoid the spectrum.

First we show that \(\operatorname{sp}(T)\) is countable. If \(m \in \mathbf{Z}^+\), then \(\{\alpha \in \operatorname{sp}(T) : |\alpha| \ge 1/m\}\) is a finite set by 10.93. Because

\begin{equation*} \operatorname{sp}(T) \subseteq \{0\} \cup \bigcup_{m=1}^{\infty} \Bigl\{\alpha \in \operatorname{sp}(T) : |\alpha| \ge \tfrac{1}{m}\Bigr\}, \end{equation*}

we see that \(\operatorname{sp}(T)\) is a countable union of finite sets together with one more point, and hence is countable.

Now fix \(n \in \mathbf{Z}^+\) and consider

\begin{equation*} A_n = \Bigl\{\alpha \in \mathbf{F} : 0 < |\alpha| < \tfrac{1}{n}\Bigr\} . \end{equation*}

The set \(A_n\) contains the interval \((0, 1/n)\) of real numbers (which lies in \(\mathbf{F}\), whether \(\mathbf{F} = \mathbf{R}\) or \(\mathbf{F} = \mathbf{C}\)) and hence is uncountable. Because \(\operatorname{sp}(T)\) is countable, \(A_n \setminus \operatorname{sp}(T) \ne \varnothing\). Choose \(\alpha_n \in A_n\) with \(\alpha_n \notin \operatorname{sp}(T)\), and set

\begin{equation*} S_n = T - \alpha_n I . \end{equation*}

Because \(\alpha_n \notin \operatorname{sp}(T)\), by the definition of the spectrum the operator \(T - \alpha_n I\) is invertible; thus each \(S_n\) is an invertible operator on \(V\). Finally,

\begin{equation*} \|T - S_n\| = \|\alpha_n I\| = |\alpha_n| \, \|I\| \le |\alpha_n| < \tfrac{1}{n} , \end{equation*}

so \(\lim_{n \to \infty} \|T - S_n\| = 0\), as desired.

Note that no compact operator on an infinite-dimensional Hilbert space is itself invertible (by 10.76), so this result shows that the invertible operators are dense enough in \(\mathcal{B}(V)\) to reach every compact operator in the limit.

Problem (10C.17)

Suppose \(T\) is a bounded operator on a Hilbert space such that \(p(T)\) is compact for some nonzero polynomial \(p\) with coefficients in \(\mathbf{F}\). Prove that \(\operatorname{sp}(T)\) is a countable set.

Solution

Let \(V\) denote the Hilbert space on which \(T\) operates.

Suppose first that \(p\) is a nonzero constant polynomial, say \(p(z) = c\) with \(c \in \mathbf{F} \setminus \{0\}\). Then \(p(T) = cI\) is compact, and hence \(I = c^{-1} p(T)\) is compact because \(\mathcal{C}(V)\) is a subspace of \(\mathcal{B}(V)\) by 10.69(a). The identity operator on an infinite-dimensional Hilbert space is not compact (see the remark following 10.67), so \(V\) is finite-dimensional. On a finite-dimensional space an operator is invertible if and only if it is injective, so \(\operatorname{sp}(T)\) consists exactly of the eigenvalues of \(T\); because eigenvectors corresponding to distinct eigenvalues are linearly independent, \(T\) has at most \(\dim V\) distinct eigenvalues. Thus \(\operatorname{sp}(T)\) is finite and hence countable.

Now suppose \(p\) has degree \(d \ge 1\).

Step 1: if \(\alpha \in \operatorname{sp}(T)\), then \(p(\alpha) \in \operatorname{sp}(p(T))\). To see this, note that the polynomial \(z \mapsto p(z) - p(\alpha)\) has coefficients in \(\mathbf{F}\) and has \(\alpha\) as a root; hence by the division algorithm there is a polynomial \(r\) with coefficients in \(\mathbf{F}\) such that

\begin{equation*} p(z) - p(\alpha) = (z - \alpha) r(z) \qquad \text{for all } z \in \mathbf{F} . \end{equation*}

Because polynomials in \(T\) commute with each other, substituting \(T\) for \(z\) gives

\begin{equation*} p(T) - p(\alpha) I = (T - \alpha I) \, r(T) = r(T) \, (T - \alpha I) . \end{equation*}

Suppose \(p(T) - p(\alpha) I\) were invertible, with (bounded) inverse \(R\). Then the displayed equation gives

\begin{equation*} (T - \alpha I) \bigl( r(T) R \bigr) = I \qquad \text{and} \qquad \bigl( R \, r(T) \bigr) (T - \alpha I) = I , \end{equation*}

so the bounded operator \(T - \alpha I\) has a bounded right inverse \(A = r(T) R\) and a bounded left inverse \(B = R \, r(T)\). Any operator with both a left and a right inverse is invertible, since \(B = B \bigl( (T - \alpha I) A \bigr) = \bigl( B (T - \alpha I) \bigr) A = A\), so \(A\) is a two-sided bounded inverse of \(T - \alpha I\). This contradicts \(\alpha \in \operatorname{sp}(T)\). Hence \(p(T) - p(\alpha) I\) is not invertible, that is, \(p(\alpha) \in \operatorname{sp}(p(T))\).

Step 2: \(\operatorname{sp}(p(T))\) is countable. Indeed, \(p(T)\) is compact, so for each \(m \in \mathbf{Z}^+\) the set \(\{\beta \in \operatorname{sp}(p(T)) : |\beta| \ge 1/m\}\) is finite by 10.93. Because

\begin{equation*} \operatorname{sp}(p(T)) \subseteq \{0\} \cup \bigcup_{m=1}^{\infty} \Bigl\{\beta \in \operatorname{sp}(p(T)) : |\beta| \ge \tfrac{1}{m}\Bigr\} , \end{equation*}

the set \(\operatorname{sp}(p(T))\) is countable.

Step 3: for each \(\beta \in \mathbf{F}\), the set \(\{z \in \mathbf{F} : p(z) = \beta\}\) is the set of roots in \(\mathbf{F}\) of the polynomial \(p - \beta\), which is a nonzero polynomial of degree \(d \ge 1\); hence that set has at most \(d\) elements. By Step 1,

\begin{equation*} \operatorname{sp}(T) \subseteq \bigcup_{\beta \in \operatorname{sp}(p(T))} \{z \in \mathbf{F} : p(z) = \beta\} . \end{equation*}

By Step 2 this is a countable union of sets each having at most \(d\) elements, and hence is countable. Thus \(\operatorname{sp}(T)\) is countable.

Problem (10C.18)

Suppose \(T \in \mathcal{B}(\ell^2)\) is defined by \(T(a_1, a_2, a_3, \ldots) = (a_2, a_3, a_4, \ldots)\). Suppose also that \(\alpha \in \mathbf{F}\) and \(|\alpha| < 1\).

(a) Show that the geometric multiplicity of \(\alpha\) as an eigenvalue of \(T\) equals \(1\).

(b) Show that the algebraic multiplicity of \(\alpha\) as an eigenvalue of \(T\) equals \(\infty\).

Solution

Recall the two definitions in play: the geometric multiplicity of an eigenvalue \(\alpha\) of \(T\) is \(\dim \operatorname{null}(T - \alpha I)\) (10.81), and the algebraic multiplicity of \(\alpha\) is the dimension of \(\bigcup_{n=1}^{\infty} \operatorname{null}(T - \alpha I)^n\) (this union is a subspace because the null spaces increase with \(n\)). Write \(S = T - \alpha I\).

(a) Suppose \(f = (a_1, a_2, \ldots) \in \ell^2\) and \(Tf = \alpha f\). Comparing coordinates in \((a_2, a_3, \ldots) = (\alpha a_1, \alpha a_2, \ldots)\) gives \(a_{n+1} = \alpha a_n\) for every \(n \in \mathbf{Z}^+\), and hence by induction \(a_n = \alpha^{n-1} a_1\) for every \(n\). Thus \(f = a_1 w\), where

\begin{equation*} w = (1, \alpha, \alpha^2, \ldots) . \end{equation*}

Because \(|\alpha| < 1\) we have \(\sum_{n=0}^{\infty} |\alpha|^{2n} = \frac{1}{1 - |\alpha|^2} < \infty\), so \(w \in \ell^2\); also \(Tw = (\alpha, \alpha^2, \alpha^3, \ldots) = \alpha w\), so \(w\) is an eigenvector of \(T\) corresponding to \(\alpha\). Combining the two paragraphs, \(\operatorname{null}(T - \alpha I) = \operatorname{span}\{w\}\), and since \(w \ne 0\) this space has dimension \(1\). Hence \(\alpha\) is an eigenvalue of \(T\) with geometric multiplicity \(1\).

(b) The main point is that \(S = T - \alpha I\) is surjective. Let \(R \in \mathcal{B}(\ell^2)\) be the right shift, \(R(a_1, a_2, \ldots) = (0, a_1, a_2, \ldots)\). Then \(R\) is an isometry, so \(\|R\| = 1\) and thus \(\|R^j\| \le 1\) for every \(j\); also \(TR = I\), and therefore

\begin{equation*} T R^{j+1} = (TR) R^j = R^j \qquad \text{for every integer } j \ge 0 . \end{equation*}

Because \(\|\alpha^j R^{j+1}\| \le |\alpha|^j\) and \(\sum_{j=0}^{\infty} |\alpha|^j < \infty\), and because \(\mathcal{B}(\ell^2)\) is a Banach space, the series

\begin{equation*} C = \sum_{j=0}^{\infty} \alpha^j R^{j+1} \end{equation*}

converges in \(\mathcal{B}(\ell^2)\) to a bounded operator \(C\) on \(\ell^2\). Applying the bounded (hence continuous) operators \(T\) and \(\alpha I\) term by term,

\begin{equation*} \begin{aligned} (T - \alpha I) C &= \sum_{j=0}^{\infty} \alpha^j T R^{j+1} - \sum_{j=0}^{\infty} \alpha^{j+1} R^{j+1} \\ &= \sum_{j=0}^{\infty} \alpha^j R^j - \sum_{i=1}^{\infty} \alpha^i R^i \\ &= R^0 = I . \end{aligned} \end{equation*}

Hence \(SC = I\), so \(S\) is surjective.

Now we show that \(\dim \operatorname{null}(S^n) = n\) for every \(n \in \mathbf{Z}^+\), by induction on \(n\). The case \(n = 1\) is part (a). Suppose \(n \in \mathbf{Z}^+\) and \(\dim \operatorname{null}(S^n) = n\). Define a linear map

\begin{equation*} \Phi : \operatorname{null}(S^{n+1}) \to \operatorname{null}(S^n), \qquad \Phi f = S f ; \end{equation*}

this makes sense because \(S^n(Sf) = S^{n+1} f = 0\) for \(f \in \operatorname{null}(S^{n+1})\). The null space of \(\Phi\) is \(\operatorname{null}(S) \cap \operatorname{null}(S^{n+1}) = \operatorname{null}(S)\), which has dimension \(1\) by (a). Moreover \(\Phi\) is surjective: if \(g \in \operatorname{null}(S^n)\), then \(f = Cg\) satisfies \(Sf = SCg = g\) and \(S^{n+1} f = S^n g = 0\), so \(f \in \operatorname{null}(S^{n+1})\) and \(\Phi f = g\). Because for a linear map the dimension of the domain equals the dimension of the null space plus the dimension of the range, we get

\begin{equation*} \dim \operatorname{null}(S^{n+1}) = 1 + \dim \operatorname{null}(S^n) = n + 1 , \end{equation*}

completing the induction.

Therefore \(\bigcup_{n=1}^{\infty} \operatorname{null}(S^n)\) contains, for each \(n \in \mathbf{Z}^+\), the subspace \(\operatorname{null}(S^n)\) of dimension \(n\); hence that union is an infinite-dimensional subspace. In other words, the algebraic multiplicity of \(\alpha\) as an eigenvalue of \(T\) equals \(\infty\).

Problem (10C.19)

Prove that the geometric multiplicity of an eigenvalue of a normal operator on a Hilbert space equals the algebraic multiplicity of that eigenvalue.

Solution

Suppose \(T\) is a normal operator on a Hilbert space \(V\) and \(\alpha \in \mathbf{F}\) is an eigenvalue of \(T\). Set \(S = T - \alpha I\). The geometric multiplicity of \(\alpha\) is \(\dim \operatorname{null} S\) and the algebraic multiplicity of \(\alpha\) is \(\dim \bigcup_{n=1}^{\infty} \operatorname{null}(S^n)\). Because \(\operatorname{null}(S) \subseteq \operatorname{null}(S^2) \subseteq \cdots\), it suffices to prove that

\begin{equation*} \operatorname{null}(S^n) = \operatorname{null}(S) \qquad \text{for every } n \in \mathbf{Z}^+ , \end{equation*}

for then the union above equals \(\operatorname{null}(S)\) and the two multiplicities are the same (whether finite or infinite).

Step 1: \(S\) is normal. By 10.12(a), 10.12(b), and 10.12(c) we have \(S^* = T^* - \overline{\alpha} I\) (with \(\overline{\alpha} = \alpha\) when \(\mathbf{F} = \mathbf{R}\)). Hence

\begin{equation*} \begin{aligned} S S^* &= T T^* - \overline{\alpha} T - \alpha T^* + |\alpha|^2 I , \\ S^* S &= T^* T - \alpha T^* - \overline{\alpha} T + |\alpha|^2 I . \end{aligned} \end{equation*}

The right sides are equal because \(T T^* = T^* T\). Thus \(S S^* = S^* S\), so \(S\) is normal.

Step 2: \(\operatorname{null}(S^2) = \operatorname{null}(S)\). The inclusion \(\operatorname{null}(S) \subseteq \operatorname{null}(S^2)\) is clear. Conversely, suppose \(S^2 f = 0\) and put \(g = Sf\), so that \(Sg = 0\). If \(g \ne 0\), then \(0\) is an eigenvalue of the normal operator \(S\) with eigenvector \(g\), so 10.56 (applied to \(S\) with eigenvalue \(0\)) gives \(S^* g = 0\); if \(g = 0\), then trivially \(S^* g = 0\). In either case \(S^* S f = S^* g = 0\), and therefore

\begin{equation*} \|S f\|^2 = \langle S f, S f \rangle = \langle f, S^* S f \rangle = \langle f, 0 \rangle = 0 , \end{equation*}

using the definition of the adjoint. Hence \(Sf = 0\), that is, \(f \in \operatorname{null}(S)\).

Step 3: induction. Suppose \(n \in \mathbf{Z}^+\) and \(\operatorname{null}(S^n) = \operatorname{null}(S)\) (the case \(n = 1\) being trivial). If \(S^{n+1} f = 0\), then \(S^n (S f) = 0\), so \(Sf \in \operatorname{null}(S^n) = \operatorname{null}(S)\), which means \(S^2 f = 0\); Step 2 now gives \(S f = 0\). Thus \(\operatorname{null}(S^{n+1}) \subseteq \operatorname{null}(S)\), and the reverse inclusion is clear. Hence \(\operatorname{null}(S^{n+1}) = \operatorname{null}(S)\), completing the induction.

Therefore \(\bigcup_{n=1}^{\infty} \operatorname{null}(T - \alpha I)^n = \operatorname{null}(T - \alpha I)\), so the algebraic multiplicity of \(\alpha\) equals its geometric multiplicity.

Problem (10C.20)

Prove that every nonzero eigenvalue of a compact operator on a Hilbert space has finite algebraic multiplicity.

Solution

Suppose \(T\) is a compact operator on a Hilbert space \(V\) and \(\alpha \in \mathbf{F} \setminus \{0\}\) is an eigenvalue of \(T\). Set \(S = T - \alpha I\); we must show that \(\bigcup_{k=1}^{\infty} \operatorname{null}(S^k)\) is finite-dimensional. Note that this union is an increasing union of subspaces, hence a subspace. Set \(\operatorname{null}(S^0) = \operatorname{null} I = \{0\}\).

Step 1: \(\operatorname{null}(S^k)\) is finite-dimensional for every \(k \in \mathbf{Z}^+\). Because \(T\) and \(I\) commute, the Binomial Theorem gives

\begin{equation*} S^k = (T - \alpha I)^k = (-\alpha)^k I + K_k , \qquad \text{where } K_k = \sum_{j=1}^{k} \binom{k}{j} (-\alpha)^{k-j} T^j . \end{equation*}

Each \(T^j\) with \(j \ge 1\) is compact by 10.69(b), and \(\mathcal{C}(V)\) is a subspace of \(\mathcal{B}(V)\) by 10.69(a); hence \(K_k\) is compact. Writing \(\beta_k = -(-\alpha)^k\), which is nonzero because \(\alpha \ne 0\), we have \(S^k = K_k - \beta_k I\), so

\begin{equation*} \operatorname{null}(S^k) = \operatorname{null}(K_k - \beta_k I) , \end{equation*}

which is finite-dimensional by 10.82 applied to the compact operator \(K_k\) and the nonzero scalar \(\beta_k\).

Step 2: there exists \(n \in \mathbf{Z}^+\) with \(\operatorname{null}(S^n) = \operatorname{null}(S^{n+1})\). Suppose not. Then

\begin{equation*} \operatorname{null}(S^0) \subsetneq \operatorname{null}(S^1) \subsetneq \operatorname{null}(S^2) \subsetneq \cdots , \end{equation*}

where the inclusions hold because \(S^k f = 0\) implies \(S^{k+1} f = 0\), and they are strict by assumption (the first one because \(\alpha\) is an eigenvalue of \(T\)). Each \(\operatorname{null}(S^k)\) is a closed subspace of \(V\), being the null space of the continuous operator \(S^k\). Fix \(n \in \mathbf{Z}^+\) and pick \(h \in \operatorname{null}(S^n) \setminus \operatorname{null}(S^{n-1})\); letting \(P\) denote the orthogonal projection of \(V\) onto the closed subspace \(\operatorname{null}(S^{n-1})\), the vector \(h - Ph\) is nonzero, lies in \(\operatorname{null}(S^n)\) (because \(Ph \in \operatorname{null}(S^{n-1}) \subseteq \operatorname{null}(S^n)\)), and is orthogonal to \(\operatorname{null}(S^{n-1})\) by 8.43. Normalizing, we obtain for each \(n \in \mathbf{Z}^+\) a vector

\begin{equation*} e_n \in \operatorname{null}(S^n) \cap \bigl( \operatorname{null}(S^{n-1}) \bigr)^{\perp} \qquad \text{with } \|e_n\| = 1 . \end{equation*}

Now suppose \(j, k \in \mathbf{Z}^+\) with \(j < k\). Since \(T = S + \alpha I\),

\begin{equation*} T e_j - T e_k = \bigl( S e_j - S e_k + \alpha e_j \bigr) - \alpha e_k . \end{equation*}

The three terms in the parentheses all lie in \(\operatorname{null}(S^{k-1})\): indeed \(S^{j-1}(S e_j) = S^j e_j = 0\) and \(j - 1 \le k - 1\), so \(S e_j \in \operatorname{null}(S^{j-1}) \subseteq \operatorname{null}(S^{k-1})\); likewise \(S e_k \in \operatorname{null}(S^{k-1})\); and \(e_j \in \operatorname{null}(S^j) \subseteq \operatorname{null}(S^{k-1})\) because \(j \le k - 1\). Because \(e_k\) is orthogonal to \(\operatorname{null}(S^{k-1})\), the Pythagorean identity gives

\begin{equation*} \|T e_j - T e_k\|^2 = \|S e_j - S e_k + \alpha e_j\|^2 + |\alpha|^2 \ge |\alpha|^2 , \end{equation*}

so \(\|T e_j - T e_k\| \ge |\alpha| > 0\) whenever \(j \ne k\). Thus the sequence \(T e_1, T e_2, \ldots\) has no convergent subsequence, even though \(e_1, e_2, \ldots\) is a bounded sequence. This contradicts the compactness of \(T\) and completes Step 2.

Step 3: with \(n\) as in Step 2, \(\operatorname{null}(S^m) = \operatorname{null}(S^{m+1})\) for every integer \(m \ge n\). We induct on \(m\). The case \(m = n\) holds by the choice of \(n\) in Step 2. Suppose \(m \ge n\) and \(\operatorname{null}(S^m) = \operatorname{null}(S^{m+1})\). If \(S^{m+2} f = 0\), then \(S^{m+1}(Sf) = 0\), so \(S f \in \operatorname{null}(S^{m+1}) = \operatorname{null}(S^m)\), which gives \(S^{m+1} f = S^m (S f) = 0\). Thus \(\operatorname{null}(S^{m+2}) \subseteq \operatorname{null}(S^{m+1})\), and the reverse inclusion is clear; hence \(\operatorname{null}(S^{m+2}) = \operatorname{null}(S^{m+1})\), completing the induction. Chaining these equalities gives \(\operatorname{null}(S^m) = \operatorname{null}(S^n)\) for every integer \(m \ge n\), while \(\operatorname{null}(S^k) \subseteq \operatorname{null}(S^n)\) for \(1 \le k \le n\).

Combining the steps,

\begin{equation*} \bigcup_{k=1}^{\infty} \operatorname{null}(T - \alpha I)^k = \operatorname{null}(S^n) , \end{equation*}

which is finite-dimensional by Step 1. Thus the algebraic multiplicity of \(\alpha\) is finite.

Problem (10C.21)

Prove that if \(T\) is a compact operator on a Hilbert space and \(\alpha\) is a nonzero eigenvalue of \(T\), then the algebraic multiplicity of \(\alpha\) as an eigenvalue of \(T\) equals the algebraic multiplicity of \(\overline{\alpha}\) as an eigenvalue of \(T^*\).

Solution

Suppose \(T\) is a compact operator on a Hilbert space \(V\) and \(\alpha \in \mathbf{F} \setminus \{0\}\) is an eigenvalue of \(T\). Set

\begin{equation*} S = T - \alpha I , \qquad \text{so that} \qquad S^* = T^* - \overline{\alpha} I \end{equation*}

by 10.12(a), 10.12(b), and 10.12(c). Also by 10.12(d) we have \((AB)^* = B^* A^*\), and hence by induction \((S^k)^* = (S^*)^k\) for every \(k \in \mathbf{Z}^+\).

First note that \(T^*\) is compact by 10.73 and that \(\overline{\alpha}\) is a nonzero eigenvalue of \(T^*\): since \(\alpha\) is an eigenvalue of \(T\) we have \(\alpha \in \operatorname{sp}(T)\), hence \(\overline{\alpha} \in \operatorname{sp}(T^*)\) (because \(T - \alpha I\) is invertible if and only if its adjoint \(T^* - \overline{\alpha} I\) is invertible, by 10.19), and the Fredholm Alternative (10.85) applied to the compact operator \(T^*\) and the nonzero scalar \(\overline{\alpha}\) shows that \(\overline{\alpha}\) is an eigenvalue of \(T^*\). This is the remark following 10.89.

By Steps 2 and 3 of the solution to Exercise 10C.20, applied to the compact operator \(T\) with the nonzero scalar \(\alpha\) and applied to the compact operator \(T^*\) with the nonzero scalar \(\overline{\alpha}\), there exist \(m, m^{\prime} \in \mathbf{Z}^+\) such that

\begin{equation*} \operatorname{null}(S^k) = \operatorname{null}(S^m) \ \text{ for all } k \ge m , \qquad \operatorname{null}\bigl( (S^*)^k \bigr) = \operatorname{null}\bigl( (S^*)^{m^{\prime}} \bigr) \ \text{ for all } k \ge m^{\prime} . \end{equation*}

Let \(N = \max\{m, m^{\prime}\}\). Then

\begin{equation*} \bigcup_{k=1}^{\infty} \operatorname{null}(T - \alpha I)^k = \operatorname{null}(S^N) \qquad \text{and} \qquad \bigcup_{k=1}^{\infty} \operatorname{null}(T^* - \overline{\alpha} I)^k = \operatorname{null}\bigl( (S^*)^N \bigr) , \end{equation*}

because in each case the null spaces increase with \(k\) and are constant from index \(N\) on. Thus it suffices to prove that

\begin{equation*} \dim \operatorname{null}(S^N) = \dim \operatorname{null}\bigl( (S^*)^N \bigr) . \end{equation*}

As in Step 1 of the solution to Exercise 10C.20, the Binomial Theorem gives

\begin{equation*} S^N = (T - \alpha I)^N = K - \beta I , \qquad \text{where } \beta = -(-\alpha)^N \ne 0 \text{ and } K = \sum_{j=1}^{N} \binom{N}{j} (-\alpha)^{N-j} T^j \end{equation*}

is compact (each \(T^j\) with \(j \ge 1\) is compact by 10.69(b), and \(\mathcal{C}(V)\) is a subspace of \(\mathcal{B}(V)\) by 10.69(a)). Taking adjoints and using 10.12,

\begin{equation*} (S^*)^N = (S^N)^* = K^* - \overline{\beta} I , \end{equation*}

and \(K^*\) is compact by 10.73.

Now apply 10.91 to the compact operator \(K\) and the nonzero scalar \(\beta\):

\begin{equation*} \dim \operatorname{null}(K - \beta I) = \dim \operatorname{null}(K^* - \overline{\beta} I) . \end{equation*}

The left side equals \(\dim \operatorname{null}(S^N)\) and the right side equals \(\dim \operatorname{null}\bigl( (S^*)^N \bigr)\). Hence

\begin{equation*} \dim \bigcup_{k=1}^{\infty} \operatorname{null}(T - \alpha I)^k = \dim \bigcup_{k=1}^{\infty} \operatorname{null}(T^* - \overline{\alpha} I)^k , \end{equation*}

which says precisely that the algebraic multiplicity of \(\alpha\) as an eigenvalue of \(T\) equals the algebraic multiplicity of \(\overline{\alpha}\) as an eigenvalue of \(T^*\). (Both numbers are finite by Exercise 10C.20.)

Problem (10C.22)

Prove that if \(V\) is a separable Hilbert space, then \(\mathcal{C}(V)\), the Banach space of compact operators on \(V\), is separable.

Solution

If \(V = \{0\}\) then \(\mathcal{C}(V) = \{0\}\), which is trivially separable. So assume \(V \neq \{0\}\).

Setup. Because \(V\) is a separable Hilbert space, 8.67 provides an orthonormal basis of \(V\), and the proof of 8.67 produces a basis indexed by \(\mathbb{Z}^+\) or by \(\{1, \ldots, N\}\) for some \(N \in \mathbb{Z}^+\). Thus we have an orthonormal basis \(\{e_k\}_{k \in \Gamma}\) of \(V\) with \(\Gamma \subseteq \mathbb{Z}^+\) (so \(\Gamma\) is countable). For \(n \in \mathbb{Z}^+\) set

\begin{equation*} \Gamma_n = \{k \in \Gamma : k \le n\}, \end{equation*}

a finite set, and let \(U_n = \operatorname{span}\{e_k\}_{k \in \Gamma_n}\). Each \(U_n\) is finite-dimensional and hence is a closed subspace of \(V\), so the orthogonal projection \(P_n = P_{U_n}\) of \(V\) onto \(U_n\) is defined. Because \(\{e_k\}_{k \in \Gamma_n}\) is an orthonormal basis of \(U_n\), 8.71 gives

\begin{equation*} P_n f = \sum_{k \in \Gamma_n} \langle f, e_k \rangle e_k \qquad \text{for all } f \in V . \end{equation*}

We record four properties of these projections.

First, \(\lim_{n \to \infty} \|f - P_n f\| = 0\) for every \(f \in V\). To see this, let \(f \in V\) and \(\varepsilon > 0\). Parseval’s identity 8.63(a) states that \(f = \sum_{k \in \Gamma} \langle f, e_k \rangle e_k\), where the sum is an unordered sum; by the definition of convergence of an unordered sum, there is a finite set \(G \subseteq \Gamma\) such that \(\big\| f - \sum_{k \in F} \langle f, e_k \rangle e_k \big\| < \varepsilon\) for every finite set \(F\) with \(G \subseteq F \subseteq \Gamma\). Because \(G\) is a finite subset of \(\Gamma \subseteq \mathbb{Z}^+\), there exists \(N \in \mathbb{Z}^+\) with \(G \subseteq \Gamma_N\); then \(G \subseteq \Gamma_n \subseteq \Gamma\) for all \(n \ge N\), and hence \(\|f - P_n f\| < \varepsilon\) for all \(n \ge N\).

Second, \(P_n\) is a bounded linear map with \(\|P_n\| \le 1\), and \(\|I - P_n\| \le 1\). Linearity is 8.37(c) and the inequality \(\|P_n f\| \le \|f\|\) is 8.37(d). For the other inequality, note that \(f - P_n f\) is orthogonal to \(P_n f\) by 8.37(a) (since \(P_n f \in U_n\)), so the Pythagorean Theorem 8.9 gives \(\|f\|^2 = \|P_n f\|^2 + \|f - P_n f\|^2\), whence \(\|(I - P_n) f\| \le \|f\|\) for every \(f \in V\).

Third, \(P_n^2 = P_n\), because \(P_n f \in U_n\) for every \(f \in V\) and \(P_n u = u\) for every \(u \in U_n\) (the latter because the displayed formula expresses \(P_n u\) as the expansion of \(u\) with respect to the orthonormal basis \(\{e_k\}_{k \in \Gamma_n}\) of \(U_n\)). Consequently \((I - P_n)^2 = I - 2P_n + P_n^2 = I - P_n\).

Fourth, \(P_n\) is self-adjoint, and hence so is \(I - P_n\) (using 10.12(a), 10.12(b), and 10.12(c)). To verify that \(P_n^* = P_n\), note that if \(f, g \in V\) then the displayed formula for \(P_n\) gives

\begin{equation*} \langle P_n f, g \rangle = \sum_{k \in \Gamma_n} \langle f, e_k \rangle \langle e_k, g \rangle \quad \text{and} \quad \langle f, P_n g \rangle = \sum_{k \in \Gamma_n} \overline{\langle g, e_k \rangle} \, \langle f, e_k \rangle = \sum_{k \in \Gamma_n} \langle f, e_k \rangle \langle e_k, g \rangle , \end{equation*}

so \(\langle P_n f, g \rangle = \langle f, P_n g \rangle\) for all \(f, g \in V\), which by the definition of the adjoint means \(P_n^* = P_n\).

Finally, recall that every compact operator is bounded (10.68); thus if \(T \in \mathcal{C}(V)\) then the adjoint \(T^*\) is defined and \(T^* \in \mathcal{B}(V)\) with \(\|T^*\| = \|T\|\), by 10.11.

Step 1: if \(T \in \mathcal{C}(V)\), then \(\lim_{n \to \infty} \|T - T P_n\| = 0\).

Suppose not. Because \(\|T - T P_1\|, \|T - T P_2\|, \ldots\) is a sequence of nonnegative numbers that does not converge to \(0\), there exist \(\varepsilon > 0\) and integers \(n_1 < n_2 < \cdots\) with \(\|T - T P_{n_j}\| > \varepsilon\) for every \(j \in \mathbb{Z}^+\). Now \(T - T P_{n_j} = T(I - P_{n_j})\), and the norm of a bounded operator is the supremum of the norms of its values on the closed unit ball; hence for each \(j\) there exists \(f_j \in V\) with \(\|f_j\| \le 1\) and

\begin{equation*} \|T(I - P_{n_j}) f_j\| > \varepsilon . \end{equation*}

Put \(g_j = (I - P_{n_j}) f_j\). Then \(\|g_j\| \le \|f_j\| \le 1\), so \(g_1, g_2, \ldots\) is a bounded sequence in \(V\). Because \(T\) is compact, there is a subsequence \(g_{j_1}, g_{j_2}, \ldots\) such that \(T g_{j_1}, T g_{j_2}, \ldots\) converges, say to \(h \in V\).

We claim \(h = 0\). Let \(g \in V\). Using the definition of the adjoint and the self-adjointness of \(I - P_n\), we get for each \(i\)

\begin{equation*} \langle T g_{j_i}, g \rangle = \langle g_{j_i}, T^* g \rangle = \big\langle (I - P_{n_{j_i}}) f_{j_i}, T^* g \big\rangle = \big\langle f_{j_i}, (I - P_{n_{j_i}}) T^* g \big\rangle, \end{equation*}

and therefore, by the Cauchy–Schwarz inequality and \(\|f_{j_i}\| \le 1\),

\begin{equation*} |\langle T g_{j_i}, g \rangle| \le \big\| (I - P_{n_{j_i}}) T^* g \big\| = \big\| T^* g - P_{n_{j_i}} T^* g \big\| . \end{equation*}

The right side tends to \(0\) as \(i \to \infty\) because \(n_{j_i} \to \infty\) and \(P_n u \to u\) for every \(u \in V\). Since \(T g_{j_i} \to h\) implies \(\langle T g_{j_i}, g \rangle \to \langle h, g \rangle\), we conclude that \(\langle h, g \rangle = 0\). As \(g \in V\) was arbitrary, taking \(g = h\) gives \(h = 0\).

But \(\|T g_{j_i}\| > \varepsilon\) for every \(i\), so \(\|h\| = \lim_{i \to \infty} \|T g_{j_i}\| \ge \varepsilon > 0\), a contradiction. This proves Step 1.

Step 2: if \(T \in \mathcal{C}(V)\), then \(\lim_{n \to \infty} \|T - P_n T P_n\| = 0\).

Write

\begin{equation*} T - P_n T P_n = (T - T P_n) + (I - P_n) T P_n . \end{equation*}

(This is an identity of bounded operators: \((T - T P_n) + (I - P_n) T P_n = T - T P_n + T P_n - P_n T P_n = T - P_n T P_n\).) For the second term, \(\|(I - P_n) T P_n\| \le \|(I - P_n) T\| \, \|P_n\| \le \|(I - P_n) T\|\). By 10.12(d) and the self-adjointness of \(I - P_n\),

\begin{equation*} \big( (I - P_n) T \big)^* = T^* (I - P_n)^* = T^* (I - P_n) = T^* - T^* P_n , \end{equation*}

so by 10.11,

\begin{equation*} \|(I - P_n) T\| = \|T^* - T^* P_n\| . \end{equation*}

Now \(T^*\) is compact by 10.73, so Step 1 applied to \(T^*\) gives \(\|T^* - T^* P_n\| \to 0\). Combining this with Step 1 applied to \(T\) yields

\begin{equation*} \|T - P_n T P_n\| \le \|T - T P_n\| + \|T^* - T^* P_n\| \longrightarrow 0 , \end{equation*}

proving Step 2.

Step 3: a countable dense subset of \(\mathcal{C}(V)\).

Let \(Q = \mathbb{Q}\) if \(\mathbf{F} = \mathbb{R}\) and \(Q = \{a + bi : a, b \in \mathbb{Q}\}\) if \(\mathbf{F} = \mathbb{C}\); in either case \(Q\) is a countable dense subset of \(\mathbf{F}\). For \(n \in \mathbb{Z}^+\) and a matrix \(c = (c_{jk})_{j, k \in \Gamma_n}\) with all entries in \(Q\), define \(A_c : V \to V\) by

\begin{equation*} A_c f = \sum_{j, k \in \Gamma_n} c_{jk} \langle f, e_k \rangle e_j . \end{equation*}

Each \(A_c\) is a linear map (because \(f \mapsto \langle f, e_k \rangle\) is linear for each \(k\) and the sum is finite), and its range is contained in \(U_n = \operatorname{span}\{e_j\}_{j \in \Gamma_n}\), which is finite-dimensional. Moreover \(A_c\) is bounded: this is the content of the basic estimate proved in the next paragraph, whose proof uses nothing from this paragraph. Hence \(A_c\) is a bounded operator whose range is finite-dimensional, so \(A_c \in \mathcal{C}(V)\) by 10.67. Let

\begin{equation*} D = \{A_c : n \in \mathbb{Z}^+ \text{ and } c \text{ is a } \Gamma_n \times \Gamma_n \text{ matrix with entries in } Q\}. \end{equation*}

For each fixed \(n\) there are only countably many such matrices (finitely many entries, each from the countable set \(Q\)), so \(D\) is a countable union of countable sets and hence is countable.

The basic estimate. If \(d = (d_{jk})_{j, k \in \Gamma_n}\) is any matrix with entries in \(\mathbf{F}\) and \(A_d\) is defined by the displayed formula, then for \(f \in V\), using orthonormality of \(\{e_j\}\), the Cauchy–Schwarz inequality in \(\mathbf{F}^{\Gamma_n}\), and Bessel’s inequality 8.57,

\begin{equation*} \begin{aligned} \|A_d f\|^2 &= \sum_{j \in \Gamma_n} \Big| \sum_{k \in \Gamma_n} d_{jk} \langle f, e_k \rangle \Big|^2 \\ &\le \sum_{j \in \Gamma_n} \Big( \sum_{k \in \Gamma_n} |d_{jk}|^2 \Big) \Big( \sum_{k \in \Gamma_n} |\langle f, e_k \rangle|^2 \Big) \\ &\le \|f\|^2 \sum_{j, k \in \Gamma_n} |d_{jk}|^2 . \end{aligned} \end{equation*}

Hence

\begin{equation*} \|A_d\| \le \Big( \sum_{j, k \in \Gamma_n} |d_{jk}|^2 \Big)^{1/2} . \end{equation*}

In particular each \(A_c\) is indeed bounded, and since \(c \mapsto A_c\) is linear in \(c\), we get \(\|A_c - A_{c^{\prime}}\| \le \big( \sum_{j,k \in \Gamma_n} |c_{jk} - c^{\prime}_{jk}|^2 \big)^{1/2}\) for two matrices \(c, c^{\prime}\) of the same size.

\(D\) is dense in \(\mathcal{C}(V)\). Suppose \(T \in \mathcal{C}(V)\) and \(\varepsilon > 0\). By Step 2 there exists \(n \in \mathbb{Z}^+\) with

\begin{equation*} \|T - P_n T P_n\| < \tfrac{\varepsilon}{2} . \end{equation*}

Let \(b_{jk} = \langle T e_k, e_j \rangle\) for \(j, k \in \Gamma_n\). For \(f \in V\),

\begin{equation*} P_n T P_n f = P_n T \Big( \sum_{k \in \Gamma_n} \langle f, e_k \rangle e_k \Big) = \sum_{k \in \Gamma_n} \langle f, e_k \rangle P_n (T e_k) = \sum_{j, k \in \Gamma_n} b_{jk} \langle f, e_k \rangle e_j , \end{equation*}

so \(P_n T P_n = A_b\) in the notation above. Because \(Q\) is dense in \(\mathbf{F}\) and \(\Gamma_n\) is finite, there is a matrix \(c = (c_{jk})_{j,k \in \Gamma_n}\) with entries in \(Q\) such that

\begin{equation*} \sum_{j, k \in \Gamma_n} |b_{jk} - c_{jk}|^2 < \frac{\varepsilon^2}{4} . \end{equation*}

The basic estimate then gives \(\|P_n T P_n - A_c\| < \frac{\varepsilon}{2}\), and therefore

\begin{equation*} \|T - A_c\| \le \|T - P_n T P_n\| + \|P_n T P_n - A_c\| < \tfrac{\varepsilon}{2} + \tfrac{\varepsilon}{2} = \varepsilon . \end{equation*}

Since \(A_c \in D\), we have shown that every element of \(\mathcal{C}(V)\) lies in the closure of \(D\).

Conversely \(D \subseteq \mathcal{C}(V)\), and \(\mathcal{C}(V)\) is a closed subset of \(\mathcal{B}(V)\) by 10.69(a); hence the closure of \(D\) is contained in \(\mathcal{C}(V)\), and the closure of \(D\) within the normed vector space \(\mathcal{C}(V)\) agrees with its closure in \(\mathcal{B}(V)\). Combining the two inclusions, the closure of \(D\) equals \(\mathcal{C}(V)\). Thus \(D\) is a countable subset of the normed vector space \(\mathcal{C}(V)\) whose closure equals \(\mathcal{C}(V)\), which by the definition 8.64 means that \(\mathcal{C}(V)\) is separable.

Remark. The key point is Step 2, which is a sharper form of the approximation in Exercise 15 in this section: there the projections are attached to an orthonormal basis of the range of the particular operator \(T\), whereas here a single sequence \(P_1, P_2, \ldots\) of projections, coming from one fixed orthonormal basis of \(V\), works simultaneously for every \(T \in \mathcal{C}(V)\). That uniformity is what allows the approximating operators \(P_n T P_n\) to be recorded by finite matrices over a fixed countable index set, and hence what makes the countable parametrization of \(D\) possible.

Exercises 10D

Problem (10D.1)

Prove that if \(T\) is a compact operator on a nonzero Hilbert space, then \(\|T\|^2\) is an eigenvalue of \(T^*T\).

Solution

Suppose \(T\) is a compact operator on a nonzero Hilbert space \(V\).

First note that \(T\) is bounded, by 10.68, so \(\|T\|\) is a well-defined finite number and the adjoint \(T^*\) is a bounded operator on \(V\) with \((T^*)^* = T\) (by 10.11). Next, \(T^*T\) is a compact operator: \(T\) is compact and \(T^*\) is bounded, so \(T^*T \in \mathcal{C}(V)\) by 10.69(b). (Alternatively, \(T^*\) is compact by 10.73 and then 10.69(b) applies in either order.)

By 10.96 we have

\begin{equation*} \|T\|^2 \in \mathrm{sp}(T^*T). \end{equation*}

Case 1: \(\|T\| \neq 0\). Then \(\|T\|^2\) is a nonzero element of \(\mathbf{F}\) lying in the spectrum of the compact operator \(T^*T\). The Fredholm Alternative 10.85, applied to the compact operator \(T^*T\) and the scalar \(\alpha = \|T\|^2 \in \mathbf{F} \setminus \{0\}\), states that \(\alpha \in \mathrm{sp}(T^*T)\) if and only if \(\alpha\) is an eigenvalue of \(T^*T\). Hence \(\|T\|^2\) is an eigenvalue of \(T^*T\).

Case 2: \(\|T\| = 0\). Then \(T = 0\) and hence \(T^*T = 0\). Because \(V \neq \{0\}\), there exists \(f \in V\) with \(f \neq 0\), and

\begin{equation*} (T^*T)f = 0 = 0 \cdot f = \|T\|^2 f . \end{equation*}

Thus \(\|T\|^2 = 0\) is an eigenvalue of \(T^*T\) (this is exactly the place where the hypothesis that \(V\) is nonzero is needed, since an operator on the zero Hilbert space has no eigenvalues).

In both cases \(\|T\|^2\) is an eigenvalue of \(T^*T\), as desired.

Problem (10D.2)

Prove that if \(T\) is a self-adjoint operator on a nonzero Hilbert space \(V\), then

\begin{equation*} \|T\| = \sup\{|\langle Tf, f\rangle| : f \in V \text{ and } \|f\| = 1\}. \end{equation*}

Solution

Recall that self-adjointness is defined only for bounded operators (10.44), so \(T\) is bounded and \(\|T\| < \infty\); recall also that \(T^* = T\) means \(\langle Tf, g\rangle = \langle f, Tg\rangle\) for all \(f, g \in V\).

Because \(V \neq \{0\}\), the set \(\{f \in V : \|f\| = 1\}\) is nonempty (if \(g \in V\) with \(g \neq 0\), then \(g/\|g\|\) belongs to it), so the supremum in the exercise is taken over a nonempty set of nonnegative numbers; call it \(M\):

\begin{equation*} M = \sup\{|\langle Tf, f\rangle| : f \in V \text{ and } \|f\| = 1\}. \end{equation*}

Also, because \(V \neq \{0\}\) we have

\begin{equation*} \|T\| = \sup\{\|Tf\| : f \in V \text{ and } \|f\| = 1\}. \end{equation*}

Indeed, \(\|T\| = \sup\{\|Tf\| : \|f\| \le 1\}\) by the definition 6.43 of the norm of a bounded linear map, and the two suprema agree: the second is over a subset of the first’s index set, so it is at most \(\|T\|\); conversely if \(0 < \|f\| \le 1\) then \(\|Tf\| = \|f\| \, \|T(f/\|f\|)\| \le \sup\{\|Tg\| : \|g\| = 1\}\), while \(f = 0\) contributes \(\|Tf\| = 0\), which is also at most that supremum (the supremum is over a nonempty set of nonnegative numbers because \(V \neq \{0\}\)).

Step 1: \(M \le \|T\|\). If \(\|f\| = 1\), then the Cauchy-Schwarz inequality 8.11 gives

\begin{equation*} |\langle Tf, f\rangle| \le \|Tf\|\,\|f\| = \|Tf\| \le \|T\|. \end{equation*}

Taking the supremum over all such \(f\) gives \(M \le \|T\|\).

Step 2: \(|\langle Tg, g\rangle| \le M\|g\|^2\) for every \(g \in V\). This is clear if \(g = 0\). If \(g \neq 0\), apply the definition of \(M\) to \(g/\|g\|\), which has norm \(1\):

\begin{equation*} \frac{|\langle Tg, g\rangle|}{\|g\|^2} = \left|\left\langle T\frac{g}{\|g\|}, \frac{g}{\|g\|}\right\rangle\right| \le M . \end{equation*}

Step 3: a polarization identity. Suppose \(g, h \in V\). Expanding by additivity of the inner product in each slot,

\begin{equation*} \begin{aligned} \langle T(g+h), g+h\rangle &= \langle Tg, g\rangle + \langle Tg, h\rangle + \langle Th, g\rangle + \langle Th, h\rangle,\\ \langle T(g-h), g-h\rangle &= \langle Tg, g\rangle - \langle Tg, h\rangle - \langle Th, g\rangle + \langle Th, h\rangle. \end{aligned} \end{equation*}

Subtracting,

\begin{equation*} \langle T(g+h), g+h\rangle - \langle T(g-h), g-h\rangle = 2\big(\langle Tg, h\rangle + \langle Th, g\rangle\big). \end{equation*}

Because \(T\) is self-adjoint, \(\langle Th, g\rangle = \langle h, Tg\rangle = \overline{\langle Tg, h\rangle}\), and therefore \(\langle Tg, h\rangle + \langle Th, g\rangle = 2\operatorname{Re}\langle Tg, h\rangle\). Hence

\begin{equation*} 4\operatorname{Re}\langle Tg, h\rangle = \langle T(g+h), g+h\rangle - \langle T(g-h), g-h\rangle . \end{equation*}

(If \(\mathbf{F} = \mathbf{R}\), then \(\operatorname{Re}\) is the identity and the identity above is the usual one.)

Step 4: the estimate. Using Step 2 on each term of the right side of the displayed identity and then the parallelogram equality 8.20,

\begin{equation*} 4|\operatorname{Re}\langle Tg, h\rangle| \le M\|g+h\|^2 + M\|g-h\|^2 = 2M\big(\|g\|^2 + \|h\|^2\big). \end{equation*}

Step 5: \(\|T\| \le M\). Suppose \(g \in V\) with \(\|g\| = 1\). If \(Tg = 0\), then trivially \(\|Tg\| \le M\). If \(Tg \neq 0\), put \(h = Tg/\|Tg\|\), so \(\|h\| = 1\) and

\begin{equation*} \langle Tg, h\rangle = \frac{\langle Tg, Tg\rangle}{\|Tg\|} = \|Tg\|, \end{equation*}

which is real, so \(\operatorname{Re}\langle Tg, h\rangle = \|Tg\|\). Step 4 now gives

\begin{equation*} 4\|Tg\| \le 2M(1 + 1) = 4M, \end{equation*}

hence \(\|Tg\| \le M\). Taking the supremum over all \(g \in V\) with \(\|g\| = 1\) yields \(\|T\| \le M\).

Combining Steps 1 and 5 gives \(\|T\| = M\), as desired.

Problem (10D.3)

Suppose \(T\) is a bounded operator on a Hilbert space \(V\) and \(U\) is a closed subspace of \(V\). Prove that the following are equivalent.

(a) \(U\) is an invariant subspace for \(T\).

(b) \(U^\perp\) is an invariant subspace for \(T^*\).

(c) \(TP_U = P_U T P_U\).

Solution

Throughout, \(P_U\) denotes the orthogonal projection of \(V\) onto the closed subspace \(U\). We will use 8.45(a), which states that \(\operatorname{range} P_U = U\) and \(\operatorname{null} P_U = U^\perp\), together with the fact (immediate from the definition of \(P_U f\) as the element of \(U\) closest to \(f\)) that \(P_U g = g\) for every \(g \in U\).

(a) \(\Longrightarrow\) (b). Assume \(U\) is invariant for \(T\), and let \(h \in U^\perp\). For every \(g \in U\) we have \(Tg \in U\), so

\begin{equation*} \langle T^*h, g\rangle = \langle h, Tg\rangle = 0 , \end{equation*}

the last equality because \(h \in U^\perp\) and \(Tg \in U\). Since \(g \in U\) was arbitrary, \(T^*h \in U^\perp\). Thus \(U^\perp\) is invariant for \(T^*\).

(b) \(\Longrightarrow\) (a). The subspace \(U^\perp\) is closed (by 8.40) and \(T^*\) is a bounded operator on \(V\). So we may apply the implication (a) \(\Rightarrow\) (b) just proved, with \(T\) replaced by \(T^*\) and \(U\) replaced by \(U^\perp\): from the hypothesis that \(U^\perp\) is invariant for \(T^*\) we conclude that \((U^\perp)^\perp\) is an invariant subspace for \((T^*)^*\). Now \((T^*)^* = T\) by 10.11, and \((U^\perp)^\perp = \overline{U} = U\) by 8.41 together with the hypothesis that \(U\) is closed. Hence \(U\) is invariant for \(T\).

(a) \(\Longrightarrow\) (c). Assume \(U\) is invariant for \(T\) and let \(f \in V\). Then \(P_U f \in U\), so \(T P_U f \in U\) by invariance, and therefore \(P_U(TP_U f) = TP_U f\). Because \(f \in V\) was arbitrary, \(P_U T P_U = T P_U\).

(c) \(\Longrightarrow\) (a). Assume \(T P_U = P_U T P_U\) and let \(f \in U\). Then \(P_U f = f\), so

\begin{equation*} Tf = T P_U f = P_U T P_U f = P_U (Tf) \in \operatorname{range} P_U = U . \end{equation*}

Thus \(U\) is invariant for \(T\).

We have shown (a) \(\Leftrightarrow\) (b) and (a) \(\Leftrightarrow\) (c), so all three statements are equivalent.

Problem (10D.4)

Suppose \(T\) is a bounded operator on a Hilbert space \(V\) and \(U\) is a closed subspace of \(V\). Prove that the following are equivalent.

(a) \(U\) and \(U^\perp\) are invariant subspaces for \(T\).

(b) \(U\) and \(U^\perp\) are invariant subspaces for \(T^*\).

(c) \(TP_U = P_U T\).

Solution

As in Exercise 3, we use 8.45: \(\operatorname{range} P_U = U\), \(\operatorname{null} P_U = U^\perp\), and \(P_{U^\perp} = I - P_U\); also \(P_U g = g\) for \(g \in U\). Note that \(U^\perp\) is a closed subspace of \(V\) (by 8.40) and that \((U^\perp)^\perp = U\) (by 8.41, since \(U\) is closed).

(a) \(\Longleftrightarrow\) (b). The equivalence of (a) and (b) in Exercise 3, applied to the closed subspace \(U\), says

\begin{equation*} U \text{ is invariant for } T \iff U^\perp \text{ is invariant for } T^* . \end{equation*}

The same equivalence, applied instead to the closed subspace \(U^\perp\), says

\begin{equation*} U^\perp \text{ is invariant for } T \iff (U^\perp)^\perp = U \text{ is invariant for } T^* . \end{equation*}

Conjoining these two equivalences gives exactly (a) \(\Leftrightarrow\) (b).

(a) \(\Longrightarrow\) (c). Assume \(U\) and \(U^\perp\) are both invariant for \(T\), and let \(f \in V\). By 8.43 we may write

\begin{equation*} f = P_U f + (f - P_U f), \qquad P_U f \in U, \quad f - P_U f = P_{U^\perp} f \in U^\perp . \end{equation*}

Applying \(T\) and using both invariance hypotheses, \(T P_U f \in U\) and \(T(f - P_U f) \in U^\perp\). Now apply \(P_U\), which is the identity on \(U\) and is \(0\) on \(U^\perp = \operatorname{null} P_U\):

\begin{equation*} P_U T f = P_U\big(T P_U f\big) + P_U\big(T(f - P_U f)\big) = T P_U f + 0 = T P_U f . \end{equation*}

Since \(f \in V\) was arbitrary, \(P_U T = T P_U\).

(c) \(\Longrightarrow\) (a). Assume \(T P_U = P_U T\).

If \(f \in U\), then \(P_U f = f\), so

\begin{equation*} Tf = T P_U f = P_U T f \in \operatorname{range} P_U = U , \end{equation*}

so \(U\) is invariant for \(T\).

If \(f \in U^\perp = \operatorname{null} P_U\), then \(P_U f = 0\), so

\begin{equation*} P_U T f = T P_U f = T0 = 0 , \end{equation*}

which means \(Tf \in \operatorname{null} P_U = U^\perp\). Thus \(U^\perp\) is invariant for \(T\).

Hence (a), (b), and (c) are equivalent.

Problem (10D.5)

Suppose \(T\) is a bounded operator on a nonseparable normed vector space \(V\). Prove that \(T\) has a closed invariant subspace other than \(\{0\}\) and \(V\).

Solution

Because the vector space \(\{0\}\) is separable, our hypothesis implies \(V \neq \{0\}\); choose \(f \in V\) with \(f \neq 0\).

Let

\begin{equation*} W = \operatorname{span}\{T^n f : n = 0, 1, 2, \dots\} \end{equation*}

(here \(T^0 = I\)), the set of all finite linear combinations of the vectors \(f, Tf, T^2f, \dots\), and let \(U = \overline{W}\).

\(U\) is a closed subspace. Indeed \(W\) is a subspace of \(V\), and the closure of a subspace of a normed vector space is a subspace: if \(g, h \in \overline{W}\) and \(\alpha \in \mathbf{F}\), choose \(g_n, h_n \in W\) with \(g_n \to g\) and \(h_n \to h\); then \(g_n + \alpha h_n \in W\) and \(g_n + \alpha h_n \to g + \alpha h\) by continuity of addition and scalar multiplication, so \(g + \alpha h \in \overline{W}\).

\(U\) is invariant for \(T\). First, \(T(W) \subseteq W\): the linear map \(T\) sends the spanning vector \(T^n f\) to \(T^{n+1} f \in W\), hence sends every finite linear combination of them into \(W\). Now suppose \(g \in U\) and choose \(g_n \in W\) with \(g_n \to g\). Since \(T\) is bounded, hence continuous, \(T g_n \to Tg\); each \(T g_n\) lies in \(W \subseteq U\), and \(U\) is closed, so \(Tg \in U\). (Compare the last bullet point of Example 10.101, which records the same fact for Hilbert spaces; the proof just given uses only that \(V\) is a normed vector space and that \(T\) is continuous.)

\(U \neq \{0\}\), because \(f \in W \subseteq U\) and \(f \neq 0\).

\(U \neq V\). Let \(\mathbf{Q}_{\mathbf{F}}\) denote \(\mathbf{Q}\) if \(\mathbf{F} = \mathbf{R}\) and \(\mathbf{Q} + i\mathbf{Q}\) if \(\mathbf{F} = \mathbf{C}\); in either case \(\mathbf{Q}_{\mathbf{F}}\) is a countable dense subset of \(\mathbf{F}\). Let

\begin{equation*} D = \left\{ \sum_{n=0}^{N} c_n T^n f : N \in \{0,1,2,\dots\} \text{ and } c_0, \dots, c_N \in \mathbf{Q}_{\mathbf{F}} \right\}. \end{equation*}

For each fixed \(N\) the corresponding set of sums is the image of the countable set \((\mathbf{Q}_{\mathbf{F}})^{N+1}\), hence countable; thus \(D\), a countable union of countable sets, is countable. Also \(D \subseteq W \subseteq U\).

\(D\) is dense in \(U\). Let \(w = \sum_{n=0}^{N} a_n T^n f \in W\) and let \(\varepsilon > 0\). Put \(C = 1 + \sum_{n=0}^{N}\|T^n f\|\) and choose \(c_n \in \mathbf{Q}_{\mathbf{F}}\) with \(|a_n - c_n| < \varepsilon/C\) for each \(n\). Then \(\sum_{n=0}^{N} c_n T^n f \in D\) and

\begin{equation*} \left\| w - \sum_{n=0}^{N} c_n T^n f \right\| \le \sum_{n=0}^{N} |a_n - c_n| \, \|T^n f\| < \varepsilon . \end{equation*}

Hence \(D\) is dense in \(W\); since \(W\) is dense in \(U = \overline{W}\), the set \(D\) is dense in \(U\).

Thus \(U\) contains a countable dense subset, so \(U\) is a separable metric space. Because \(V\) is not separable, \(U \neq V\).

Therefore \(U\) is a closed invariant subspace for \(T\) with \(U \neq \{0\}\) and \(U \neq V\).

Problem (10D.6)

Suppose \(T\) is an operator on a Banach space \(V\) with dimension greater than \(2\). Prove that \(T\) has an invariant subspace other than \(\{0\}\) and \(V\).

[For this exercise, \(T\) is not assumed to be bounded and the invariant subspace is not required to be closed.]

Solution

Because \(\dim V > 2\), we have \(V \neq \{0\}\); fix \(f \in V\) with \(f \neq 0\) and set

\begin{equation*} U = \operatorname{span}\{T^n f : n = 0, 1, 2, \dots\} \end{equation*}

(the set of finite linear combinations, with \(T^0 = I\)). Then \(U\) is a subspace of \(V\), it is invariant for \(T\) (the linear map \(T\) sends each spanning vector \(T^n f\) to \(T^{n+1} f \in U\)), and \(U \neq \{0\}\) because \(f \in U\). So the whole issue is whether \(U\) can be forced to be a proper subspace. We treat two cases.

Case 1: \(V\) is infinite-dimensional.

We claim \(U \neq V\), so that \(U\) itself is the desired invariant subspace. Suppose instead \(U = V\). For \(n \in \mathbf{Z}^+\) let

\begin{equation*} V_n = \operatorname{span}\{f, Tf, \dots, T^{n-1}f\}. \end{equation*}

Every element of \(U\) is a finite linear combination of the vectors \(T^k f\), so \(V = U = \bigcup_{n=1}^{\infty} V_n\).

Each \(V_n\) is finite-dimensional (its dimension is at most \(n\)), and every finite-dimensional subspace of a normed vector space is closed (Exercise 9 in Section 6D). Each \(V_n\) is a proper subspace of \(V\), because \(V\) is infinite-dimensional while \(\dim V_n \le n\).

A proper subspace \(W\) of a normed vector space \(V\) has empty interior. To see this, suppose some open ball \(B(g, r)\) with \(r > 0\) satisfies \(B(g,r) \subseteq W\); note \(g \in W\). If \(h \in V\) with \(h \neq 0\), then \(\big\|\big(g + \tfrac{r}{2\|h\|}h\big) - g\big\| = r/2 < r\), so \(g + \tfrac{r}{2\|h\|}h \in W\); subtracting \(g \in W\) and multiplying by \(2\|h\|/r\) shows \(h \in W\). Hence \(W = V\), contradicting properness.

Thus \(V\) is the countable union of the closed sets \(V_1, V_2, \dots\), each with empty interior. But \(V\) is a Banach space and hence a complete metric space, so Baire’s Theorem 6.76(a) says that \(V\) is not the countable union of closed subsets with empty interior. This contradiction shows \(U \neq V\), completing Case 1.

Case 2: \(V\) is finite-dimensional, say \(\dim V = n\) with \(n \ge 3\).

If \(U \neq V\) we are done, so assume \(U = V\).

The list \(f, Tf, T^2f, \dots\) spans \(V\), and \(\dim V = n\). Let \(m\) be the smallest positive integer for which the list \(f, Tf, \dots, T^{m}f\) is linearly dependent; such \(m\) exists and \(m \le n\), because any \(n+1\) vectors in \(V\) are linearly dependent. Then \(f, Tf, \dots, T^{m-1}f\) is linearly independent and \(T^m f \in \operatorname{span}\{f, \dots, T^{m-1}f\}\). Write \(W = \operatorname{span}\{f, \dots, T^{m-1}f\}\). Then \(T^k f \in W\) for every \(k \ge 0\), by induction on \(k\): this is immediate for \(k \le m-1\), and if \(T^k f \in W\) for some \(k \ge m-1\), then \(T^{k+1} f \in T(W) \subseteq \operatorname{span}\{Tf, \dots, T^{m}f\} \subseteq W\), the last inclusion because \(T^m f \in W\). Since the vectors \(T^k f\) span \(U\), we get

\begin{equation*} V = U = \operatorname{span}\{f, Tf, \dots, T^{m-1}f\}. \end{equation*}

Hence \(m = n\) and \(f, Tf, \dots, T^{n-1}f\) is a basis of \(V\). Write

\begin{equation*} T^n f = -\big(a_0 f + a_1 Tf + \dots + a_{n-1}T^{n-1}f\big) \end{equation*}

for some \(a_0, \dots, a_{n-1} \in \mathbf{F}\), and define the monic polynomial

\begin{equation*} p(z) = z^n + a_{n-1}z^{n-1} + \dots + a_1 z + a_0 , \end{equation*}

so that \(p(T)f = 0\).

In fact \(p(T) = 0\) as an operator on \(V\): for each \(k \ge 0\), polynomials in \(T\) commute with \(T\), so \(p(T)T^k f = T^k p(T) f = 0\); since the vectors \(T^k f\) span \(V\) and \(p(T)\) is linear, \(p(T) = 0\).

Factor \(p\) into monic irreducible factors over \(\mathbf{F}\):

\begin{equation*} p = q_1 q_2 \cdots q_r . \end{equation*}

If \(\mathbf{F} = \mathbf{C}\), every irreducible monic polynomial is linear; if \(\mathbf{F} = \mathbf{R}\), every irreducible monic polynomial has degree \(1\) or \(2\) (the standard factorization of a real polynomial into linear factors and irreducible quadratic factors). In either case \(d := \deg q_1 \le 2\). Because \(n \ge 3\), we have \(1 \le d \le 2 \le n - 1\). Also \(r \ge 2\), because otherwise \(p = q_1\) would be irreducible of degree \(n \ge 3\), which is impossible over \(\mathbf{R}\) or \(\mathbf{C}\). Set \(s := q_2 \cdots q_r\), a monic polynomial with \(\deg s = n - d\); since \(1 \le d \le 2 < n\), we have \(1 \le n - d \le n - 1\).

Here is the key observation, used twice: if \(q\) is a monic polynomial with \(1 \le \deg q \le n-1\), then \(q(T)f \neq 0\). Indeed, writing \(\deg q = e \le n-1\), the vector \(q(T)f\) is a linear combination of \(f, Tf, \dots, T^{e}f\) in which the coefficient of \(T^{e}f\) equals \(1\); since \(e \le n-1\) and \(f, Tf, \dots, T^{n-1}f\) is linearly independent, this combination is not \(0\).

Now let

\begin{equation*} N = \operatorname{null} q_1(T). \end{equation*}

\(N\) is invariant for \(T\): if \(q_1(T)g = 0\), then \(q_1(T)(Tg) = T\big(q_1(T)g\big) = 0\), so \(Tg \in N\).

\(N \neq V\): by the key observation applied to \(q_1\) (which is monic with \(1 \le d \le n-1\)), \(q_1(T)f \neq 0\), so \(f \notin N\).

\(N \neq \{0\}\): by the key observation applied to \(s\) (monic with \(1 \le n-d \le n-1\)), the vector \(g := s(T)f\) is nonzero; and

\begin{equation*} q_1(T)g = q_1(T)s(T)f = p(T)f = 0 , \end{equation*}

so \(0 \neq g \in N\).

Thus \(N\) is an invariant subspace for \(T\) other than \(\{0\}\) and \(V\), completing Case 2.

In both cases \(T\) has an invariant subspace other than \(\{0\}\) and \(V\), as desired. (The hypothesis \(\dim V > 2\) cannot be dropped: counterclockwise rotation by a right angle on \(\mathbf{R}^2\) has no invariant subspace other than \(\{0\}\) and \(\mathbf{R}^2\). The completeness hypothesis was used only in Case 1, to rule out an infinite-dimensional normed vector space spanned by countably many vectors.)

Problem (10D.7)

Suppose \(T\) is a self-adjoint compact operator on a Hilbert space that has only finitely many distinct eigenvalues. Prove that \(T\) has finite-dimensional range.

Solution

Let \(V\) denote the Hilbert space on which \(T\) operates.

By the Spectral Theorem for self-adjoint compact operators 10.106(b), there exist a countable set \(\Omega\), an orthonormal family \(\{e_k\}_{k \in \Omega}\) in \(V\), and a family \(\{\alpha_k\}_{k \in \Omega}\) in \(\mathbf{R} \setminus \{0\}\) such that

\begin{equation*} Tf = \sum_{k \in \Omega} \alpha_k \langle f, e_k\rangle e_k \end{equation*}

for every \(f \in V\). We show that \(\Omega\) is a finite set; the conclusion follows at once from that.

Step 1: each \(e_j\) with \(j \in \Omega\) is an eigenvector of \(T\) corresponding to the eigenvalue \(\alpha_j\). Indeed, taking \(f = e_j\) in the displayed formula and using orthonormality (so that \(\langle e_j, e_k\rangle\) equals \(1\) if \(k = j\) and equals \(0\) otherwise), every term of the sum vanishes except the term with \(k = j\), giving

\begin{equation*} Te_j = \alpha_j e_j . \end{equation*}

Since \(\|e_j\| = 1\), we have \(e_j \neq 0\), so \(\alpha_j\) is indeed an eigenvalue of \(T\); moreover \(\alpha_j \neq 0\).

Step 2: only finitely many distinct values occur among \(\{\alpha_k\}_{k \in \Omega}\). By Step 1 each \(\alpha_k\) is an eigenvalue of \(T\), and by hypothesis \(T\) has only finitely many distinct eigenvalues. So there are distinct nonzero real numbers \(\lambda_1, \dots, \lambda_p\) (with \(p \ge 0\)) such that \(\{\alpha_k : k \in \Omega\} \subseteq \{\lambda_1, \dots, \lambda_p\}\).

Step 3: \(\Omega\) is finite. For each \(j \in \{1, \dots, p\}\) let

\begin{equation*} \Omega_j = \{k \in \Omega : \alpha_k = \lambda_j\} . \end{equation*}

By Step 1, \(\{e_k\}_{k \in \Omega_j}\) is an orthonormal family contained in \(\operatorname{null}(T - \lambda_j I)\). Because \(T\) is compact and \(\lambda_j \neq 0\), the subspace \(\operatorname{null}(T - \lambda_j I)\) is finite-dimensional by 10.82. An orthonormal family is linearly independent (if \(\sum_{k \in F} c_k e_k = 0\) for a finite \(F \subseteq \Omega_j\), taking the inner product with \(e_i\) gives \(c_i = 0\) for each \(i \in F\)), so

\begin{equation*} \#\Omega_j \le \dim \operatorname{null}(T - \lambda_j I) < \infty . \end{equation*}

By Step 2, \(\Omega = \Omega_1 \cup \dots \cup \Omega_p\), a union of finitely many finite sets. Hence \(\Omega\) is finite.

Step 4: conclusion. Because \(\Omega\) is finite, the formula from 10.106(b) exhibits \(Tf\), for each \(f \in V\), as a finite linear combination of \(e_k\) for \(k \in \Omega\). Therefore

\begin{equation*} \operatorname{range} T \subseteq \operatorname{span}\{e_k : k \in \Omega\}, \end{equation*}

and the subspace on the right has dimension at most \(\#\Omega < \infty\). Hence \(\operatorname{range} T\) is finite-dimensional, with

\begin{equation*} \dim \operatorname{range} T \le \#\Omega . \end{equation*}

(If \(\Omega = \emptyset\), the formula gives \(T = 0\) and \(\operatorname{range} T = \{0\}\), which is finite-dimensional.)

Problem (10D.8)

(a) Prove that if \(T\) is a self-adjoint compact operator on a Hilbert space, then there exists a self-adjoint compact operator \(S\) such that \(S^3 = T\).

(b) Prove that if \(T\) is a normal compact operator on a complex Hilbert space, then there exists a normal compact operator \(S\) such that \(S^2 = T\).

Solution

Both parts are obtained by applying a function to the eigenvalues in a spectral decomposition. We first isolate the analytic content in a lemma about diagonal operators, which will be used in both parts.

Lemma. Suppose \(\{e_k\}_{k \in \Omega}\) is an orthonormal family in a Hilbert space \(V\), where \(\Omega\) is a countable set, and suppose \(\{\beta_k\}_{k \in \Omega}\) is a family in \(\mathbb{F}\) such that

\begin{equation*} \{k \in \Omega : |\beta_k| \geq \delta\} \text{ is finite for every } \delta > 0. \end{equation*}

Then the formula

\begin{equation*} S f = \sum_{k \in \Omega} \beta_k \langle f, e_k \rangle e_k \end{equation*}

defines a compact operator \(S\) on \(V\) with \(\|S\| \leq \sup_{k \in \Omega} |\beta_k|\), and

  • (i) \(S^* f = \sum_{k \in \Omega} \overline{\beta_k} \langle f, e_k \rangle e_k\) for all \(f \in V\);
  • (ii) \(S^n f = \sum_{k \in \Omega} \beta_k^{\,n} \langle f, e_k \rangle e_k\) for all \(f \in V\) and all \(n \in \mathbb{Z}^+\);
  • (iii) \(S e_j = \beta_j e_j\) for each \(j \in \Omega\), and \(S g = 0\) for every \(g \in \bigl(\overline{\operatorname{span}}\{e_k\}_{k \in \Omega}\bigr)^\perp\).

Proof of Lemma. If \(\Omega = \emptyset\) then \(S = 0\) and every assertion is trivial, so assume \(\Omega \neq \emptyset\). The hypothesis on \(\{\beta_k\}\) applied with \(\delta = 1\) shows that all but finitely many \(\beta_k\) satisfy \(|\beta_k| < 1\); hence \(M := \sup_{k \in \Omega} |\beta_k| < \infty\). If \(f \in V\), then Bessel’s inequality (8.57) gives

\begin{equation*} \sum_{k \in \Omega} |\beta_k \langle f, e_k \rangle|^2 \leq M^2 \sum_{k \in \Omega} |\langle f, e_k \rangle|^2 \leq M^2 \|f\|^2 < \infty, \end{equation*}

so the unordered sum defining \(S f\) converges by 8.54(a), and 8.54(b) gives \(\|S f\|^2 = \sum_{k \in \Omega} |\beta_k|^2 |\langle f, e_k \rangle|^2 \leq M^2 \|f\|^2\). Linearity of \(S\) is clear from the linearity of \(f \mapsto \langle f, e_k \rangle\) and the fact that unordered sums are additive. Thus \(S\) is a bounded operator with \(\|S\| \leq M\).

Taking the inner product of the defining sum with a fixed \(e_j\) and using continuity of the inner product together with orthonormality gives

\begin{equation*} \langle S f, e_j \rangle = \beta_j \langle f, e_j \rangle \quad \text{for each } j \in \Omega. \end{equation*}

Statement (iii) follows by taking \(f = e_j\) (all terms with \(k \neq j\) vanish) and by noting that \(\langle g, e_k \rangle = 0\) for all \(k\) when \(g\) is orthogonal to every \(e_k\). Statement (ii) follows by induction: if \(S^n f = \sum_k \beta_k^{\,n} \langle f, e_k \rangle e_k\), then applying the displayed identity to \(S^n f\) gives \(\langle S^n f, e_j \rangle = \beta_j^{\,n} \langle f, e_j\rangle\), hence

\begin{equation*} S^{n+1} f = S(S^n f) = \sum_{k \in \Omega} \beta_k \langle S^n f, e_k \rangle e_k = \sum_{k \in \Omega} \beta_k^{\,n+1} \langle f, e_k \rangle e_k . \end{equation*}

For (i), let \(S^{\prime}\) denote the operator defined by the family \(\{\overline{\beta_k}\}_{k \in \Omega}\), which satisfies the same hypotheses. For \(f, g \in V\), expanding both inner products through the (convergent) unordered sums gives

\begin{equation*} \langle S f, g \rangle = \sum_{k \in \Omega} \beta_k \langle f, e_k \rangle \langle e_k, g \rangle = \sum_{k \in \Omega} \langle f, e_k \rangle \overline{\overline{\beta_k} \langle g, e_k \rangle} = \langle f, S^{\prime} g \rangle . \end{equation*}

Because this holds for all \(f, g \in V\), we have \(S^* = S^{\prime}\), proving (i).

Finally we prove that \(S\) is compact. If \(\Omega\) is finite, then \(\operatorname{range} S\) is finite-dimensional and \(S\) is compact by 10.67. So suppose \(\Omega\) is infinite and enumerate it as \(\Omega = \{k_1, k_2, \ldots\}\) with the \(k_j\) distinct. The hypothesis says that for each \(\delta > 0\) only finitely many \(j\) satisfy \(|\beta_{k_j}| \geq \delta\); hence \(\lim_{j \to \infty} \beta_{k_j} = 0\). For \(n \in \mathbb{Z}^+\) define

\begin{equation*} S_n f = \sum_{j = 1}^{n} \beta_{k_j} \langle f, e_{k_j} \rangle e_{k_j} . \end{equation*}

Each \(S_n\) is a bounded operator with finite-dimensional range and hence is compact by 10.67. The operator \(S - S_n\) is the diagonal operator determined by the family that agrees with \(\{\beta_k\}\) on \(\{k_{n+1}, k_{n+2}, \ldots\}\) and is \(0\) elsewhere, so by the norm estimate already proved,

\begin{equation*} \|S - S_n\| \leq \sup_{j > n} |\beta_{k_j}| \longrightarrow 0 \quad \text{as } n \to \infty . \end{equation*}

Because \(\mathcal{C}(V)\) is a closed subspace of \(\mathcal{B}(V)\) [10.69(a)], we conclude that \(S\) is compact. This completes the proof of the lemma.

Proof of (a). Suppose \(T\) is a self-adjoint compact operator on a Hilbert space \(V\). By 10.106(b) there exist a countable set \(\Omega\), an orthonormal family \(\{e_k\}_{k \in \Omega}\) in \(V\), and a family \(\{\alpha_k\}_{k \in \Omega}\) in \(\mathbb{R} \setminus \{0\}\) such that

\begin{equation*} T f = \sum_{k \in \Omega} \alpha_k \langle f, e_k \rangle e_k \end{equation*}

for every \(f \in V\). Taking \(f = e_j\) shows that \(T e_j = \alpha_j e_j\), so each \(\alpha_k\) is a nonzero eigenvalue of \(T\).

We claim that \(\{k \in \Omega : |\alpha_k| \geq \delta\}\) is finite for every \(\delta > 0\). Fix \(\delta > 0\). Every eigenvalue of \(T\) belongs to \(\operatorname{sp}(T)\), so by 10.93 the set \(A_\delta\) of eigenvalues \(\alpha\) of \(T\) with \(|\alpha| \geq \delta\) is finite. For each \(\alpha \in A_\delta\), the vectors \(\{e_k : \alpha_k = \alpha\}\) form an orthonormal, hence linearly independent, subset of \(\operatorname{null}(T - \alpha I)\), which is finite-dimensional by 10.82 (as \(\alpha \neq 0\)); therefore \(\{k \in \Omega : \alpha_k = \alpha\}\) is finite. Since \(\{k : |\alpha_k| \geq \delta\}\) is the union over the finite set \(A_\delta\) of these finite sets, it is finite, proving the claim.

For \(k \in \Omega\) let \(\beta_k = \alpha_k^{1/3}\) denote the real cube root of \(\alpha_k\). Then \(|\beta_k| = |\alpha_k|^{1/3}\), so \(|\beta_k| \geq \delta\) if and only if \(|\alpha_k| \geq \delta^3\); hence \(\{\beta_k\}\) satisfies the hypothesis of the lemma. Let \(S\) be the operator produced by the lemma from \(\{\beta_k\}\). Then \(S\) is compact, and \(S\) is self-adjoint by (i) of the lemma because each \(\beta_k\) is real. Finally, (ii) of the lemma gives

\begin{equation*} S^3 f = \sum_{k \in \Omega} \beta_k^{\,3} \langle f, e_k \rangle e_k = \sum_{k \in \Omega} \alpha_k \langle f, e_k \rangle e_k = T f \end{equation*}

for every \(f \in V\). Thus \(S^3 = T\), as desired.

Proof of (b). Suppose \(T\) is a normal compact operator on a complex Hilbert space \(V\). By the Spectral Theorem for normal compact operators (10.107) there is an orthonormal basis \(\{e_k\}_{k \in \Gamma}\) of \(V\) consisting of eigenvectors of \(T\); say \(T e_k = \alpha_k e_k\) with \(\alpha_k \in \mathbb{C}\). If \(f \in V\), then Parseval’s identity 8.63(a) gives \(f = \sum_{k \in \Gamma} \langle f, e_k \rangle e_k\), and applying the bounded operator \(T\) (which is continuous, hence commutes with the convergent unordered sum) yields

\begin{equation*} T f = \sum_{k \in \Gamma} \alpha_k \langle f, e_k \rangle e_k . \end{equation*}

Let \(\Omega = \{k \in \Gamma : \alpha_k \neq 0\}\); the terms with \(k \notin \Omega\) vanish, so \(T f = \sum_{k \in \Omega} \alpha_k \langle f, e_k \rangle e_k\) for all \(f \in V\). Exactly as in the proof of (a), each \(\alpha_k\) with \(k \in \Omega\) is a nonzero eigenvalue of \(T\), and 10.93 together with 10.82 shows that \(\{k \in \Omega : |\alpha_k| \geq \delta\}\) is finite for every \(\delta > 0\); in particular \(\Omega\) is countable.

For each \(k \in \Omega\) choose \(\beta_k \in \mathbb{C}\) with \(\beta_k^{\,2} = \alpha_k\) (every complex number has a square root). Then \(|\beta_k| = |\alpha_k|^{1/2}\), so \(|\beta_k| \geq \delta\) if and only if \(|\alpha_k| \geq \delta^2\), and hence \(\{\beta_k\}_{k \in \Omega}\) satisfies the hypothesis of the lemma. Let \(S\) be the corresponding operator; the lemma shows that \(S\) is compact and, by (ii),

\begin{equation*} S^2 f = \sum_{k \in \Omega} \beta_k^{\,2} \langle f, e_k \rangle e_k = \sum_{k \in \Omega} \alpha_k \langle f, e_k \rangle e_k = T f \end{equation*}

for every \(f \in V\), so \(S^2 = T\).

It remains to check that \(S\) is normal. By (iii) of the lemma, \(S e_j = \beta_j e_j\) for \(j \in \Omega\) and \(S e_j = 0\) for \(j \in \Gamma \setminus \Omega\) (such an \(e_j\) is orthogonal to every \(e_k\) with \(k \in \Omega\)). Thus every element of the orthonormal basis \(\{e_k\}_{k \in \Gamma}\) of \(V\) is an eigenvector of \(S\). Because \(S\) is a bounded operator on a complex Hilbert space admitting an orthonormal basis of eigenvectors, 10.103(b) implies that \(S\) is normal, completing the proof.

Problem (10D.9)

Suppose \(T\) is a compact normal operator on a nonzero Hilbert space \(V\). Prove that there is a subspace of \(V\) with dimension \(1\) or \(2\) that is an invariant subspace for \(T\).

[If \(\mathbb{F} = \mathbb{C}\), the desired result follows immediately from the Spectral Theorem for compact normal operators. Thus you can assume that \(\mathbb{F} = \mathbb{R}\).]

Solution

Suppose first that \(\mathbb{F} = \mathbb{C}\). By the Spectral Theorem for normal compact operators (10.107), there is an orthonormal basis of \(V\) consisting of eigenvectors of \(T\). Because \(V \neq \{0\}\), this orthonormal basis is nonempty; let \(f\) be one of its elements. Then \(f\) is an eigenvector of \(T\), so \(\operatorname{span}\{f\}\) is a \(1\)-dimensional invariant subspace for \(T\) (see the second bullet point of 10.101).

Now suppose \(\mathbb{F} = \mathbb{R}\). Define

\begin{equation*} A = \frac{T + T^*}{2}, \qquad B = \frac{T - T^*}{2}, \end{equation*}

so that \(T = A + B\). Because \(T\) is compact, so is \(T^*\) (by 10.73); because \(\mathcal{C}(V)\) is a subspace of \(\mathcal{B}(V)\) [10.69(a)], both \(A\) and \(B\) are compact. Clearly \(A^* = A\) and \(B^* = -B\).

Since \(T\) is normal we have \(T T^* = T^* T\), and therefore

\begin{equation*} TA = \tfrac{1}{2}(T^2 + T T^*) = \tfrac{1}{2}(T^2 + T^* T) = AT, \qquad BA = \tfrac{1}{2}(TA - T^*A) = \tfrac{1}{2}(AT - AT^*) = AB, \end{equation*}

where the middle equality in the second chain uses \(TA = AT\) and the analogous computation \(T^* A = \tfrac12(T^*T + (T^*)^2) = \tfrac12(TT^* + (T^*)^2) = A T^*\). Thus \(A\) commutes with both \(T\) and \(B\).

The operator \(A\) is a self-adjoint compact operator on the nonzero Hilbert space \(V\), so by 10.99 either \(\|A\|\) or \(-\|A\|\) is an eigenvalue of \(A\). Fix such an eigenvalue \(\alpha \in \mathbb{R}\) and set

\begin{equation*} U = \operatorname{null}(A - \alpha I), \end{equation*}

so \(U \neq \{0\}\). Since \(A\) is bounded, \(U\) is a closed subspace of \(V\), hence itself a nonzero Hilbert space. If \(f \in U\), then

\begin{equation*} (A - \alpha I)(B f) = B(A - \alpha I) f = B(0) = 0, \end{equation*}

using \(AB = BA\); hence \(B f \in U\). Thus \(U\) is an invariant subspace for \(B\). The identical computation with \(T\) in place of \(B\) (using \(AT = TA\)) shows that \(U\) is an invariant subspace for \(T\) as well.

Let \(C = B|_U\), regarded as an operator on \(U\). Because \(U\) is a closed invariant subspace for the compact operator \(B\), the operator \(C\) is compact on \(U\) (as noted in the discussion preceding 10.102). For \(f, g \in U\) we have \(Bf, Bg \in U\) and

\begin{equation*} \langle C f, g \rangle = \langle B f, g \rangle = \langle f, B^* g \rangle = -\langle f, B g \rangle = \langle f, -C g \rangle, \end{equation*}

so \(C^* = -C\). Consequently \((C^2)^* = (C^*)^2 = (-C)^2 = C^2\), so \(C^2\) is a self-adjoint operator on \(U\), and it is compact by 10.69(b). Moreover, for every \(g \in U\),

\begin{equation*} \langle C^2 g, g \rangle = \langle C g, C^* g \rangle = -\langle C g, C g \rangle = -\|C g\|^2 \leq 0 . \end{equation*}

Because \(U \neq \{0\}\), applying 10.99 to the self-adjoint compact operator \(C^2\) on \(U\) produces an eigenvalue \(\lambda\) of \(C^2\); choose \(g \in U\) with \(g \neq 0\) and \(C^2 g = \lambda g\). The displayed inequality gives \(\lambda \|g\|^2 = \langle C^2 g, g\rangle \leq 0\), so \(\lambda \leq 0\).

Note that for every \(f \in U\) we have \(A f = \alpha f\), hence

\begin{equation*} T f = A f + B f = \alpha f + C f . \end{equation*}

Case 1: \(\lambda = 0\). Then \(\|C g\|^2 = -\langle C^2 g, g \rangle = 0\), so \(C g = 0\) and therefore \(T g = \alpha g\). Thus \(\operatorname{span}\{g\}\) is a \(1\)-dimensional subspace of \(V\) that is invariant for \(T\).

Case 2: \(\lambda < 0\). Let \(W = \operatorname{span}\{g, C g\}\), a subspace of \(U\) and hence of \(V\). First, \(\dim W = 2\): if \(g\) and \(Cg\) were linearly dependent, then (since \(g \neq 0\)) we would have \(C g = c g\) for some \(c \in \mathbb{R}\), giving \(\lambda g = C^2 g = c^2 g\) and hence \(c^2 = \lambda < 0\), which is impossible for a real number \(c\). Second, \(W\) is invariant for \(T\): using \(T f = \alpha f + C f\) for \(f \in U\),

\begin{equation*} T g = \alpha g + C g \in W, \qquad T(C g) = \alpha\, C g + C^2 g = \lambda g + \alpha\, C g \in W, \end{equation*}

and every element of \(W\) is a linear combination of \(g\) and \(Cg\). Thus \(W\) is a \(2\)-dimensional invariant subspace for \(T\).

In either case we have produced a subspace of \(V\) of dimension \(1\) or \(2\) that is invariant for \(T\), completing the proof.

Problem (10D.10)

Suppose \(T\) is a self-adjoint compact operator on a Hilbert space and \(\|T\| \leq \frac{1}{4}\). Prove that there exists a self-adjoint compact operator \(S\) such that \(S^2 + S = T\).

Solution

The idea is to solve the scalar equation \(s^2 + s = t\) for each eigenvalue \(t\) of \(T\) and then assemble the solutions into an operator. The hypothesis \(\|T\| \leq \frac14\) guarantees that a real solution exists for each eigenvalue.

Step 1: a scalar function. For \(t \in [-\tfrac14, \tfrac14]\) we have \(1 + 4t \geq 0\), so we may define

\begin{equation*} \varphi(t) = \frac{-1 + \sqrt{1 + 4t}}{2}. \end{equation*}

Writing \(u = \sqrt{1+4t}\), so that \(u^2 = 1 + 4t\) and \(\varphi(t) = \frac{u-1}{2}\), we compute

\begin{equation*} \varphi(t)^2 + \varphi(t) = \frac{u^2 - 2u + 1}{4} + \frac{2u - 2}{4} = \frac{u^2 - 1}{4} = t . \end{equation*}

Furthermore, rationalizing the numerator gives

\begin{equation*} \varphi(t) = \frac{(\sqrt{1+4t} - 1)(\sqrt{1+4t}+1)}{2(\sqrt{1+4t}+1)} = \frac{2t}{1 + \sqrt{1+4t}}, \end{equation*}

and since \(1 + \sqrt{1+4t} \geq 1\) for \(t \in [-\tfrac14, \tfrac14]\), we obtain the bound

\begin{equation*} |\varphi(t)| \leq 2 |t| \qquad \text{for all } t \in [-\tfrac14, \tfrac14]. \end{equation*}

Step 2: the spectral decomposition of \(T\). Let \(V\) denote the Hilbert space on which \(T\) acts. If \(V = \{0\}\) then \(S = 0\) works, so assume \(V \neq \{0\}\). By 10.106(b) there exist a countable set \(\Omega\), an orthonormal family \(\{e_k\}_{k \in \Omega}\) in \(V\), and a family \(\{\alpha_k\}_{k \in \Omega}\) in \(\mathbb{R} \setminus \{0\}\) such that

\begin{equation*} T f = \sum_{k \in \Omega} \alpha_k \langle f, e_k \rangle e_k \end{equation*}

for every \(f \in V\). Taking \(f = e_j\) shows \(T e_j = \alpha_j e_j\), so each \(\alpha_k\) is a nonzero eigenvalue of \(T\); in particular \(|\alpha_k| \, \|e_k\| = \|T e_k\| \leq \|T\| \, \|e_k\|\), so

\begin{equation*} |\alpha_k| \leq \|T\| \leq \tfrac14 \qquad \text{for every } k \in \Omega . \end{equation*}

We also claim that \(\{k \in \Omega : |\alpha_k| \geq \delta\}\) is finite for each \(\delta > 0\). Indeed, eigenvalues of \(T\) lie in \(\operatorname{sp}(T)\), so by 10.93 only finitely many distinct values \(\alpha\) with \(|\alpha| \geq \delta\) occur among the \(\alpha_k\); and for each such nonzero \(\alpha\), the vectors \(\{e_k : \alpha_k = \alpha\}\) are orthonormal, hence linearly independent, in the finite-dimensional space \(\operatorname{null}(T - \alpha I)\) [finite-dimensional by 10.82], so only finitely many \(k\) have \(\alpha_k = \alpha\). A finite union of finite sets is finite, proving the claim.

Step 3: the diagonal operator lemma. Suppose \(\{\beta_k\}_{k \in \Omega}\) is a family in \(\mathbb{R}\) such that \(\{k \in \Omega : |\beta_k| \geq \delta\}\) is finite for every \(\delta > 0\). We claim that

\begin{equation*} S f = \sum_{k \in \Omega} \beta_k \langle f, e_k \rangle e_k \end{equation*}

defines a self-adjoint compact operator \(S\) on \(V\) satisfying \(\|S\| \leq M := \sup_{k \in \Omega} |\beta_k|\) and \(S^2 f = \sum_{k \in \Omega} \beta_k^{\,2} \langle f, e_k \rangle e_k\).

Applying the hypothesis with \(\delta = 1\) shows \(M < \infty\). For \(f \in V\), Bessel’s inequality (8.57) gives \(\sum_{k} |\beta_k \langle f, e_k\rangle|^2 \leq M^2 \|f\|^2 < \infty\), so the unordered sum converges by 8.54(a), and \(\|S f\|^2 = \sum_k |\beta_k|^2 |\langle f, e_k\rangle|^2 \leq M^2 \|f\|^2\) by 8.54(b). Linearity is clear, so \(S\) is a bounded operator with \(\|S\| \leq M\). Taking the inner product with \(e_j\) and using orthonormality and continuity of the inner product gives \(\langle S f, e_j \rangle = \beta_j \langle f, e_j\rangle\); applying \(S\) once more therefore yields \(S^2 f = \sum_k \beta_k \langle S f, e_k\rangle e_k = \sum_k \beta_k^{\,2} \langle f, e_k \rangle e_k\). For self-adjointness, expand both sides through the convergent sums: for \(f, g \in V\),

\begin{equation*} \langle S f, g \rangle = \sum_{k \in \Omega} \beta_k \langle f, e_k \rangle \langle e_k, g \rangle = \sum_{k \in \Omega} \langle f, e_k\rangle \overline{\beta_k \langle g, e_k \rangle} = \langle f, S g \rangle, \end{equation*}

where we used that each \(\beta_k\) is real; hence \(S^* = S\). For compactness: if \(\Omega\) is finite then \(\operatorname{range} S\) is finite-dimensional and \(S\) is compact by 10.67. If \(\Omega\) is infinite, enumerate \(\Omega = \{k_1, k_2, \ldots\}\) with distinct \(k_j\); the hypothesis forces \(\lim_{j \to \infty} \beta_{k_j} = 0\). Setting \(S_n f = \sum_{j=1}^n \beta_{k_j} \langle f, e_{k_j}\rangle e_{k_j}\), each \(S_n\) is compact by 10.67, and \(S - S_n\) is the diagonal operator for the family that equals \(\beta_k\) on \(\{k_{n+1}, k_{n+2}, \ldots\}\) and \(0\) elsewhere, so the norm bound gives \(\|S - S_n\| \leq \sup_{j > n} |\beta_{k_j}| \to 0\). Because \(\mathcal{C}(V)\) is closed in \(\mathcal{B}(V)\) [10.69(a)], \(S\) is compact. This proves the claim.

Step 4: conclusion. For \(k \in \Omega\) set \(\beta_k = \varphi(\alpha_k)\), which is defined because \(|\alpha_k| \leq \frac14\) by Step 2, and is real. By Step 1, \(|\beta_k| \leq 2|\alpha_k|\), so

\begin{equation*} \{k \in \Omega : |\beta_k| \geq \delta\} \subseteq \{k \in \Omega : |\alpha_k| \geq \tfrac{\delta}{2}\}, \end{equation*}

which is finite for every \(\delta > 0\) by Step 2. Hence the family \(\{\beta_k\}_{k \in \Omega}\) satisfies the hypothesis of Step 3; let \(S\) be the resulting self-adjoint compact operator. For every \(f \in V\), Step 3 and Step 1 give

\begin{equation*} (S^2 + S) f = \sum_{k \in \Omega} \bigl(\beta_k^{\,2} + \beta_k\bigr) \langle f, e_k \rangle e_k = \sum_{k \in \Omega} \bigl(\varphi(\alpha_k)^2 + \varphi(\alpha_k)\bigr) \langle f, e_k \rangle e_k = \sum_{k \in \Omega} \alpha_k \langle f, e_k \rangle e_k = T f . \end{equation*}

Thus \(S\) is a self-adjoint compact operator with \(S^2 + S = T\), as desired.

Problem (10D.11)

For \(k \in \mathbb{Z}\), define \(g_k \in L^2\bigl((-\pi, \pi]\bigr)\) and \(h_k \in L^2\bigl((-\pi, \pi]\bigr)\) by

\begin{equation*} g_k(t) = \frac{1}{\sqrt{2\pi}}\, e^{it/2} e^{ikt} \qquad \text{and} \qquad h_k(t) = \frac{1}{\sqrt{2\pi}}\, e^{ikt}; \end{equation*}

here we are assuming that \(\mathbb{F} = \mathbb{C}\).

(a) Use the conclusion of Example 10.108 to show that \(\{g_k\}_{k \in \mathbb{Z}}\) is an orthonormal basis of \(L^2\bigl((-\pi, \pi]\bigr)\).

(b) Use (a) to show that \(\{h_k\}_{k \in \mathbb{Z}}\) is an orthonormal basis of \(L^2\bigl((-\pi, \pi]\bigr)\).

(c) Use (b) to show that the orthonormal family in the third bullet point of Example 8.51 is an orthonormal basis of \(L^2\bigl((-\pi, \pi]\bigr)\).

Solution

Throughout, all \(L^2\) spaces are taken with respect to Lebesgue measure and \(\mathbb{F} = \mathbb{C}\). Recall the conclusion of Example 10.108: if for \(k \in \mathbb{Z}\) we define \(u_k \in L^2([0,1])\) by

\begin{equation*} u_k(x) = e^{i(2k+1)\pi x}, \end{equation*}

then \(\{u_k\}_{k \in \mathbb{Z}}\) is an orthonormal basis of \(L^2([0,1])\).

We will use twice the elementary observation that a surjective linear isometry between Hilbert spaces carries orthonormal bases to orthonormal bases. Precisely, suppose \(\Psi : V \to W\) is a linear bijection between Hilbert spaces with \(\|\Psi f\| = \|f\|\) for all \(f \in V\). Then \(\Psi\) preserves inner products, since by the polarization identity the inner product is determined by the norm; hence \(\Psi\) carries an orthonormal family to an orthonormal family. Moreover \(\Psi\) and \(\Psi^{-1}\) are continuous, so \(\Psi\) is a homeomorphism, and therefore \(\Psi\bigl(\overline{\operatorname{span}}\, A\bigr) = \overline{\operatorname{span}}\, \Psi(A)\) for every subset \(A\) of \(V\). Consequently if \(\{f_k\}_{k \in \Gamma}\) is an orthonormal basis of \(V\), then \(\{\Psi f_k\}_{k \in \Gamma}\) is an orthonormal family whose closed span is \(\Psi(V) = W\), i.e. an orthonormal basis of \(W\).

We also record the change-of-variables fact used below: if \(a > 0\), \(b \in \mathbb{R}\), and \(\varphi\) is a nonnegative Borel measurable function on \(\mathbb{R}\), then

\begin{equation*} \int \varphi(a x + b) \, dx = \frac{1}{a} \int \varphi(y) \, dy . \end{equation*}

This holds for \(\varphi = \chi_E\) because Lebesgue measure is translation invariant (2.7) and satisfies \(|cE| = c\,|E|\) for \(c > 0\) (immediate from the definition of outer measure, since the lengths of the intervals in a cover get multiplied by \(c\)); it then extends to nonnegative simple functions by linearity and to all nonnegative measurable \(\varphi\) by the Monotone Convergence Theorem (3.11).

Part (a). Define \(\Phi : L^2([0,1]) \to L^2\bigl((-\pi,\pi]\bigr)\) by

\begin{equation*} (\Phi f)(t) = \frac{1}{\sqrt{2\pi}}\, f\Bigl( \frac{t}{2\pi} + \frac{1}{2} \Bigr) \qquad \text{for } t \in (-\pi, \pi] . \end{equation*}

As \(t\) runs over \((-\pi, \pi]\), the point \(x = \frac{t}{2\pi} + \frac12\) runs over \((0, 1]\), so the formula makes sense. The function \(t \mapsto f\bigl(\frac{t}{2\pi} + \frac12\bigr)\) is Borel measurable, being the composition of the Borel measurable \(f\) with a continuous map. Moreover \(\Phi\) is well defined on equivalence classes: the change-of-variables fact with \(\varphi = \chi_E\) shows that the affine map \(t \mapsto \frac{t}{2\pi} + \frac12\) pulls sets of measure \(0\) back to sets of measure \(0\), so replacing \(f\) by a function equal to it almost everywhere changes \(\Phi f\) only on a set of measure \(0\). Now extend \(|f|^2\) by \(0\) to all of \(\mathbb{R}\) and apply the change-of-variables fact to this extension with \(a = \frac{1}{2\pi}\) and \(b = \frac12\); because the extension vanishes off \([0,1]\), both integrals below are over the indicated sets, and we get

\begin{equation*} \|\Phi f\|^2 = \int_{-\pi}^{\pi} \frac{1}{2\pi} \Bigl| f\Bigl(\frac{t}{2\pi} + \frac12\Bigr) \Bigr|^2 \, dt = \frac{1}{2\pi} \cdot 2\pi \int_0^1 |f(x)|^2 \, dx = \|f\|^2 . \end{equation*}

Thus \(\Phi\) is a linear isometry. It is surjective, because the formula \((\Phi^{-1} v)(x) = \sqrt{2\pi}\, v(2\pi x - \pi)\) defines a linear map \(L^2\bigl((-\pi,\pi]\bigr) \to L^2([0,1])\) (by the same change of variables) that is a two-sided inverse of \(\Phi\).

Now compute \(\Phi u_k\). For \(t \in (-\pi, \pi]\),

\begin{equation*} (\Phi u_k)(t) = \frac{1}{\sqrt{2\pi}}\, e^{i(2k+1)\pi \left( \frac{t}{2\pi} + \frac12 \right)} = \frac{1}{\sqrt{2\pi}}\, e^{i(2k+1)t/2}\, e^{i(2k+1)\pi/2} = c_k\, g_k(t), \end{equation*}

where \(c_k = e^{i(2k+1)\pi/2}\) satisfies \(|c_k| = 1\), and where we used \(e^{i(2k+1)t/2} = e^{it/2} e^{ikt}\).

By the first paragraph, \(\{\Phi u_k\}_{k \in \mathbb{Z}}\) is an orthonormal basis of \(L^2\bigl((-\pi,\pi]\bigr)\). Since \(g_k = \overline{c_k}\, \Phi u_k\) with \(|\overline{c_k}| = 1\), the family \(\{g_k\}_{k \in \mathbb{Z}}\) is again orthonormal (because \(\langle g_j, g_k\rangle = \overline{c_j} c_k \langle \Phi u_j, \Phi u_k\rangle\) and \(\overline{c_k} c_k = 1\)), and \(\operatorname{span}\{g_k\}_{k \in \mathbb{Z}} = \operatorname{span}\{\Phi u_k\}_{k \in \mathbb{Z}}\) because each family is obtained from the other by multiplying by nonzero scalars. Hence the two families have the same closed span, namely \(L^2\bigl((-\pi,\pi]\bigr)\). Therefore \(\{g_k\}_{k \in \mathbb{Z}}\) is an orthonormal basis of \(L^2\bigl((-\pi,\pi]\bigr)\).

Part (b). Let \(m(t) = e^{-it/2}\) for \(t \in (-\pi,\pi]\) and define \(M : L^2\bigl((-\pi,\pi]\bigr) \to L^2\bigl((-\pi,\pi]\bigr)\) by \(M v = m v\). Because \(|m(t)| = 1\) for every \(t\), we have \(|(Mv)(t)| = |v(t)|\) pointwise and hence \(\|M v\| = \|v\|\); so \(M\) is a linear isometry. It is bijective, with inverse given by multiplication by \(\overline{m}(t) = e^{it/2}\). By the observation in the first paragraph, \(M\) carries orthonormal bases to orthonormal bases. Finally,

\begin{equation*} (M g_k)(t) = e^{-it/2} \cdot \frac{1}{\sqrt{2\pi}} e^{it/2} e^{ikt} = \frac{1}{\sqrt{2\pi}} e^{ikt} = h_k(t) . \end{equation*}

Thus \(\{h_k\}_{k \in \mathbb{Z}} = \{M g_k\}_{k \in \mathbb{Z}}\) is an orthonormal basis of \(L^2\bigl((-\pi,\pi]\bigr)\) by (a).

Part (c). The third bullet point of Example 8.51 defines, for \(k \in \mathbb{Z}\), the function \(e_k : (-\pi, \pi] \to \mathbb{R}\) by

\begin{equation*} e_k(t) = \begin{cases} \frac{1}{\sqrt{\pi}} \sin(kt) & \text{if } k > 0, \\[2pt] \frac{1}{\sqrt{2\pi}} & \text{if } k = 0, \\[2pt] \frac{1}{\sqrt{\pi}} \cos(kt) & \text{if } k < 0. \end{cases} \end{equation*}

We express these in terms of the \(h_k\). First, \(e_0 = h_0\). Next, for \(n \in \mathbb{Z}^+\) we have

\begin{equation*} (h_n - h_{-n})(t) = \frac{e^{int} - e^{-int}}{\sqrt{2\pi}} = \frac{2i \sin(nt)}{\sqrt{2\pi}}, \qquad (h_n + h_{-n})(t) = \frac{e^{int} + e^{-int}}{\sqrt{2\pi}} = \frac{2\cos(nt)}{\sqrt{2\pi}} , \end{equation*}

and therefore, using \(\frac{\sqrt{2\pi}}{2\sqrt{\pi}} = \frac{1}{\sqrt 2}\) together with \(e_n(t) = \frac{1}{\sqrt{\pi}} \sin(nt)\) and \(e_{-n}(t) = \frac{1}{\sqrt{\pi}}\cos(-nt) = \frac{1}{\sqrt{\pi}}\cos(nt)\), we obtain the identities in \(L^2\bigl((-\pi,\pi]\bigr)\)

\begin{equation*} e_n = \frac{1}{i\sqrt{2}}\, (h_n - h_{-n}), \qquad e_{-n} = \frac{1}{\sqrt{2}}\, (h_n + h_{-n}) . \end{equation*}

Solving these two equations for \(h_n\) and \(h_{-n}\) gives, for \(n \in \mathbb{Z}^+\),

\begin{equation*} h_n = \frac{1}{\sqrt{2}}\bigl( e_{-n} + i\, e_n \bigr), \qquad h_{-n} = \frac{1}{\sqrt{2}}\bigl( e_{-n} - i\, e_n \bigr) . \end{equation*}

First we verify that \(\{e_k\}_{k \in \mathbb{Z}}\) is orthonormal, which now follows from the orthonormality of \(\{h_k\}_{k \in \mathbb{Z}}\) established in (b). Indeed \(\|e_0\| = \|h_0\| = 1\), and for \(n \in \mathbb{Z}^+\) the Pythagorean Theorem (8.9) gives

\begin{equation*} \|e_n\|^2 = \tfrac{1}{2}\bigl( \|h_n\|^2 + \|h_{-n}\|^2 \bigr) = 1, \qquad \|e_{-n}\|^2 = \tfrac{1}{2}\bigl( \|h_n\|^2 + \|h_{-n}\|^2 \bigr) = 1 , \end{equation*}

while

\begin{equation*} \langle e_n, e_{-n} \rangle = \frac{1}{i\sqrt 2} \cdot \overline{\Bigl(\frac{1}{\sqrt2}\Bigr)} \, \bigl\langle h_n - h_{-n},\, h_n + h_{-n} \bigr\rangle = \frac{1}{2i}\bigl( 1 - 1 \bigr) = 0 . \end{equation*}

If \(m, n \in \mathbb{Z}^+\) with \(m \neq n\), then \(e_{\pm n}\) is a linear combination of \(h_n, h_{-n}\) and \(e_{\pm m}\) is a linear combination of \(h_m, h_{-m}\), and these four indices are distinct, so all such inner products vanish; similarly \(e_0 = h_0\) is orthogonal to \(e_{\pm n}\) for every \(n \in \mathbb{Z}^+\). Hence \(\{e_k\}_{k \in \mathbb{Z}}\) is an orthonormal family, as asserted in 8.51.

Finally, the displayed formulas show that every \(e_k\) lies in \(\operatorname{span}\{h_j\}_{j \in \mathbb{Z}}\) and every \(h_k\) lies in \(\operatorname{span}\{e_j\}_{j \in \mathbb{Z}}\). Therefore

\begin{equation*} \operatorname{span}\{e_k\}_{k \in \mathbb{Z}} = \operatorname{span}\{h_k\}_{k \in \mathbb{Z}}, \end{equation*}

and taking closures and using (b) gives

\begin{equation*} \overline{\operatorname{span}}\,\{e_k\}_{k \in \mathbb{Z}} = \overline{\operatorname{span}}\,\{h_k\}_{k \in \mathbb{Z}} = L^2\bigl((-\pi,\pi]\bigr) . \end{equation*}

Thus the orthonormal family \(\{e_k\}_{k \in \mathbb{Z}}\) of the third bullet point of Example 8.51 is an orthonormal basis of \(L^2\bigl((-\pi,\pi]\bigr)\), as desired.

Problem (10D.12)

Suppose \(T\) is a compact operator on a Hilbert space. Prove that \(s_1(T) = \|T\|\).

Solution

Let \(V\) denote the Hilbert space on which \(T\) acts. Note that \(T^*\) is compact (10.73) and hence \(T^* T\) is compact by 10.69(b); moreover \((T^* T)^* = T^* T^{**} = T^* T\), so \(T^* T\) is self-adjoint.

We first record an upper bound. Suppose \(\alpha\) is an eigenvalue of \(T^* T\), say \(T^* T f = \alpha f\) with \(f \in V\) and \(f \neq 0\). Then

\begin{equation*} \alpha \|f\|^2 = \langle \alpha f, f \rangle = \langle T^* T f, f \rangle = \langle T f, T f \rangle = \|T f\|^2 \leq \|T\|^2 \|f\|^2 . \end{equation*}

Dividing by \(\|f\|^2 > 0\) gives

\begin{equation*} \alpha \leq \|T\|^2 \qquad \text{for every eigenvalue } \alpha \text{ of } T^* T . \end{equation*}

(The same computation shows \(\alpha \geq 0\), which is the observation used at the start of the proof of 10.113.)

Next we record a lower bound. Suppose \(V \neq \{0\}\) and \(\|T\| \neq 0\). By 10.96 applied to \(T\), we have \(\|T\|^2 \in \operatorname{sp}(T^* T)\). Since \(\|T\|^2 \neq 0\) and \(T^* T\) is compact, the Fredholm Alternative (10.85) implies that

\begin{equation*} \|T\|^2 \text{ is an eigenvalue of } T^* T . \end{equation*}

Because \(\|T\|^2 > 0\), this is a positive eigenvalue of \(T^*T\).

Now we consider two cases.

Case 1: \(T^* T\) has no positive eigenvalue. Then by the second bullet point of Definition 10.116, \(s_n(T) = 0\) for every \(n \in \mathbb{Z}^+\); in particular \(s_1(T) = 0\). On the other hand, if we had \(\|T\| \neq 0\), then \(V \neq \{0\}\) (since \(T = 0\) on the zero space) and the lower bound above would produce a positive eigenvalue of \(T^* T\), contradicting the case assumption. Hence \(\|T\| = 0 = s_1(T)\).

Case 2: \(T^* T\) has at least one positive eigenvalue. Let \(P\) denote the set of positive eigenvalues of \(T^* T\) and fix \(\alpha_0 \in P\). Every element of \(P\) lies in \(\operatorname{sp}(T^*T)\), so by 10.93 the set \(\{\alpha \in P : \alpha \geq \alpha_0\}\) is finite; it is also nonempty (it contains \(\alpha_0\)), so it has a largest element, and that element is \(\max P\), since every element of \(P\) not in this set is smaller than \(\alpha_0\). By the first bullet point of Definition 10.116, the singular values of \(T\) are the positive square roots of the elements of \(P\), listed with multiplicity in decreasing order; the largest of them is therefore

\begin{equation*} s_1(T) = \sqrt{\max P}. \end{equation*}

Because \(T^* T \neq 0\) in this case, we have \(T \neq 0\) and \(V \neq \{0\}\), so the lower bound applies: \(\|T\|^2 \in P\), whence \(\|T\|^2 \leq \max P = s_1(T)^2\). The upper bound applied to \(\alpha = \max P\) gives \(s_1(T)^2 = \max P \leq \|T\|^2\). Combining the two inequalities yields \(s_1(T)^2 = \|T\|^2\), and since \(s_1(T) \geq 0\) and \(\|T\| \geq 0\) we conclude \(s_1(T) = \|T\|\).

In both cases \(s_1(T) = \|T\|\), as desired.

Problem (10D.13)

Suppose \(T\) is a compact operator on a Hilbert space and \(n \in \mathbb{Z}^+\). Prove that \(\dim \operatorname{range} T < n\) if and only if \(s_n(T) = 0\).

Solution

Let \(V\) denote the Hilbert space on which \(T\) acts. The operator \(T^* T\) is compact (by 10.73 and 10.69(b)) and self-adjoint, and as shown at the start of the proof of 10.113 all of its eigenvalues are nonnegative.

Step 1: a singular value decomposition. Applying 10.106(b) to \(T^* T\) gives a countable set \(\Omega\), an orthonormal family \(\{e_k\}_{k \in \Omega}\) in \(V\), and a family \(\{\alpha_k\}_{k \in \Omega}\) of nonzero real numbers, necessarily positive by the previous paragraph, such that

\begin{equation*} (T^* T) f = \sum_{k \in \Omega} \alpha_k \langle f, e_k \rangle e_k \end{equation*}

for all \(f \in V\). Put \(s_k = \sqrt{\alpha_k} > 0\) and \(h_k = \frac{T e_k}{s_k}\) for \(k \in \Omega\). As shown in the proof of 10.113, \(\{h_k\}_{k \in \Omega}\) is an orthonormal family in \(V\) and

\begin{equation*} T f = \sum_{k \in \Omega} s_k \langle f, e_k \rangle h_k \end{equation*}

for all \(f \in V\).

Step 2: \(\dim \operatorname{range} T\) and the cardinality of \(\Omega\). Because \(T e_k = s_k h_k\) with \(s_k \neq 0\), each \(h_k\) belongs to \(\operatorname{range} T\). The family \(\{h_k\}_{k \in \Omega}\) is orthonormal and hence linearly independent, so

\begin{equation*} \dim \operatorname{range} T \geq |\Omega| . \end{equation*}

Conversely, if \(\Omega\) is finite then the displayed formula for \(T f\) shows \(\operatorname{range} T \subseteq \operatorname{span}\{h_k\}_{k \in \Omega}\), so \(\dim \operatorname{range} T \leq |\Omega|\). Therefore

\begin{equation*} \dim \operatorname{range} T = |\Omega| , \end{equation*}

where both sides are \(\infty\) when \(\Omega\) is infinite.

Step 3: the number of positive singular values equals \(|\Omega|\). We claim that for every \(\alpha > 0\),

\begin{equation*} \operatorname{null}(T^* T - \alpha I) = \operatorname{span}\{ e_k : k \in \Omega, \ \alpha_k = \alpha \} . \end{equation*}

The inclusion \(\supseteq\) is immediate, since taking \(f = e_j\) in the formula for \(T^*T\) gives \((T^* T) e_j = \alpha_j e_j\). For \(\subseteq\), suppose \((T^* T) f = \alpha f\). Taking the inner product of the formula for \((T^*T)f\) with \(e_j\) gives \(\langle (T^*T) f, e_j\rangle = \alpha_j \langle f, e_j\rangle\), while \(\langle \alpha f, e_j\rangle = \alpha \langle f, e_j \rangle\); hence \((\alpha_j - \alpha)\langle f, e_j\rangle = 0\), so \(\langle f, e_j \rangle = 0\) whenever \(\alpha_j \neq \alpha\). Therefore

\begin{equation*} \alpha f = (T^* T) f = \sum_{k \in \Omega,\ \alpha_k = \alpha} \alpha \langle f, e_k \rangle e_k , \end{equation*}

and dividing by \(\alpha \neq 0\) shows that \(f = \sum_{\alpha_k = \alpha} \langle f, e_k \rangle e_k\). The index set \(\{k \in \Omega : \alpha_k = \alpha\}\) is finite, because the \(e_k\) with \(\alpha_k = \alpha\) form an orthonormal subset of the finite-dimensional space \(\operatorname{null}(T^*T - \alpha I)\) [finite-dimensional by 10.82, as \(\alpha \neq 0\)]; so the sum is a genuine finite linear combination and \(f \in \operatorname{span}\{e_k : \alpha_k = \alpha\}\). This proves the claim.

Two consequences follow. First, the positive eigenvalues of \(T^* T\) are exactly the numbers occurring in \(\{\alpha_k\}_{k \in \Omega}\): if \(\alpha > 0\) were an eigenvalue not occurring among the \(\alpha_k\), the claim would give \(\operatorname{null}(T^* T - \alpha I) = \{0\}\), a contradiction. Second, for each positive eigenvalue \(\alpha\) of \(T^* T\), the geometric multiplicity \(\dim \operatorname{null}(T^* T - \alpha I)\) equals the number of \(k \in \Omega\) with \(\alpha_k = \alpha\), since the \(e_k\) with \(\alpha_k = \alpha\) form an orthonormal (hence linearly independent) spanning set of that space.

By Definition 10.116, the singular values \(s_1(T) \geq s_2(T) \geq \cdots\) list the positive square roots of the positive eigenvalues of \(T^* T\) in decreasing order, each such square root \(s\) repeated as many times as the geometric multiplicity of \(s^2\) as an eigenvalue of \(T^* T\), and then padded with zeros if that list is finite. By the two consequences just established, that list of positive numbers is exactly a rearrangement of the family \(\{s_k\}_{k \in \Omega}\) into decreasing order. Hence, writing \(N \in \{0, 1, 2, \ldots\} \cup \{\infty\}\) for the number of positive singular values of \(T\),

\begin{equation*} N = |\Omega|, \qquad \text{and} \qquad s_m(T) > 0 \iff m \leq N \quad \text{for } m \in \mathbb{Z}^+ . \end{equation*}

Step 4: conclusion. Fix \(n \in \mathbb{Z}^+\). Using Steps 2 and 3,

\begin{equation*} s_n(T) = 0 \iff \text{not } (n \leq N) \iff N < n \iff |\Omega| < n \iff \dim \operatorname{range} T < n , \end{equation*}

which is the desired equivalence.

Problem (10D.14)

Suppose \(T\) is a compact operator on a Hilbert space \(V\) with singular value decomposition

\begin{equation*} T f = \sum_{k=1}^{\infty} s_k(T) \langle f, e_k \rangle h_k \end{equation*}

for all \(f \in V\). For \(n \in \mathbb{Z}^+\), define \(T_n : V \to V\) by

\begin{equation*} T_n f = \sum_{k=1}^{n} s_k(T) \langle f, e_k \rangle h_k . \end{equation*}

Prove that \(\lim_{n \to \infty} \|T - T_n\| = 0\).

[This exercise gives another proof, in addition to the proof suggested by Exercise 15 in Section 10C, that an operator on a Hilbert space is compact if and only if it is the limit of bounded operators with finite-dimensional range.]

Solution

Here \(\{e_k\}_{k \in \mathbb{Z}^+}\) and \(\{h_k\}_{k \in \mathbb{Z}^+}\) are orthonormal families in \(V\), as in 10.113, and \(s_1(T) \geq s_2(T) \geq \cdots\) are the singular values of \(T\) as defined in 10.116. We do not need to assume that every \(s_k(T)\) is positive; only the decreasing ordering is used below.

Step 1: the estimate \(\|T - T_n\| \leq s_{n+1}(T)\). Fix \(n \in \mathbb{Z}^+\) and \(f \in V\). Since \(T_n f\) is the sum of the first \(n\) terms of the convergent series defining \(T f\),

\begin{equation*} (T - T_n) f = \sum_{k = n+1}^{\infty} s_k(T) \langle f, e_k \rangle h_k . \end{equation*}

The family \(\{h_k\}_{k > n}\) is orthonormal, so 8.54(b) gives

\begin{equation*} \|(T - T_n) f\|^2 = \sum_{k = n+1}^{\infty} s_k(T)^2 \, |\langle f, e_k \rangle|^2 . \end{equation*}

Because the singular values are arranged in decreasing order, \(s_k(T) \leq s_{n+1}(T)\) for every \(k \geq n+1\). Hence, using Bessel’s inequality (8.57) for the orthonormal family \(\{e_k\}_{k \in \mathbb{Z}^+}\),

\begin{equation*} \|(T - T_n) f\|^2 \leq s_{n+1}(T)^2 \sum_{k = n+1}^{\infty} |\langle f, e_k \rangle|^2 \leq s_{n+1}(T)^2 \, \|f\|^2 . \end{equation*}

Taking the supremum over \(f \in V\) with \(\|f\| \leq 1\) gives

\begin{equation*} \|T - T_n\| \leq s_{n+1}(T) . \end{equation*}

(In fact equality holds, since \((T - T_n) e_{n+1} = s_{n+1}(T) h_{n+1}\) has norm \(s_{n+1}(T)\) and \(\|e_{n+1}\| = 1\), but we only need the inequality.)

Step 2: \(\lim_{n \to \infty} s_n(T) = 0\). The operator \(T^* T\) is compact (by 10.73 and 10.69(b)) and self-adjoint. If \(T^* T\) has only finitely many positive eigenvalues, then by the second bullet point of Definition 10.116 we have \(s_k(T) = 0\) for all sufficiently large \(k\), and the limit is \(0\); so assume from now on that \(T^* T\) has infinitely many positive eigenvalues, in which case every \(s_k(T)\) is given by the first bullet point of 10.116. Let \(\delta > 0\). By that bullet point, each \(s_k(T)\) is the positive square root of a positive eigenvalue of \(T^* T\), and each such square root \(s\) is listed exactly \(\dim \operatorname{null}(T^* T - s^2 I)\) times. Now:

  • The eigenvalues of \(T^* T\) lie in \(\operatorname{sp}(T^* T)\), so by 10.93 the set of eigenvalues \(\alpha\) of \(T^* T\) with \(\alpha \geq \delta^2\) is finite.
  • For each such \(\alpha\), the space \(\operatorname{null}(T^* T - \alpha I)\) is finite-dimensional by 10.82, since \(\alpha \neq 0\).

Therefore the value \(\sqrt{\alpha}\) is listed only finitely many times for each of the finitely many eigenvalues \(\alpha \geq \delta^2\), and so

\begin{equation*} \{ k \in \mathbb{Z}^+ : s_k(T) \geq \delta \} \end{equation*}

is a finite set. Consequently there exists \(N \in \mathbb{Z}^+\) such that \(s_k(T) < \delta\) for all \(k > N\). As \(\delta > 0\) was arbitrary, \(\lim_{n \to \infty} s_n(T) = 0\).

Step 3: conclusion. Combining Steps 1 and 2,

\begin{equation*} 0 \leq \|T - T_n\| \leq s_{n+1}(T) \longrightarrow 0 \qquad \text{as } n \to \infty, \end{equation*}

so \(\lim_{n \to \infty} \|T - T_n\| = 0\), as desired.

As the bracketed remark indicates, this yields the stated characterization of compact operators. Each \(T_n\) is a bounded operator whose range is contained in \(\operatorname{span}\{h_1, \ldots, h_n\}\) and hence is finite-dimensional; thus every compact operator with a singular value decomposition as above is a limit in norm of bounded operators with finite-dimensional range. (If the index set of the singular value decomposition is finite, then \(T\) itself has finite-dimensional range and there is nothing to prove.) Conversely, bounded operators with finite-dimensional range are compact by 10.67, and \(\mathcal{C}(V)\) is closed in \(\mathcal{B}(V)\) by 10.69(a), so any norm limit of such operators is compact.

Problem (10D.15)

Suppose \(T\) is a compact operator on a Hilbert space \(V\) and \(n \in \mathbb{Z}^+\). Prove that

\begin{equation*} \inf\{\|T - S\| : S \in \mathcal{B}(V) \text{ and } \dim \operatorname{range} S < n\} = s_n(T). \end{equation*}

Solution

Throughout, fix a singular value decomposition of \(T\), which exists by 10.113: there are a countable set \(\Omega\), orthonormal families \(\{e_k\}_{k \in \Omega}\) and \(\{h_k\}_{k \in \Omega}\) in \(V\), and positive numbers \(\{s_k\}_{k \in \Omega}\) such that

\begin{equation*} Tf = \sum_{k \in \Omega} s_k \langle f, e_k \rangle h_k \end{equation*}

for every \(f \in V\). As explained in the discussion preceding 10.116, the family \(\{s_k\}_{k \in \Omega}\) is uniquely determined by \(T\): its members are the positive square roots of the positive eigenvalues of \(T^*T\), each repeated as many times as the geometric multiplicity of the corresponding eigenvalue, and this multiplicity is finite by 10.82. Hence \(\Omega\) is countable and the numbers \(s_k\) may be arranged in decreasing order (if there are infinitely many of them they form a sequence with limit \(0\), by 10.93). Concretely, let \(\Omega^{\prime}\) be \(\{1, \ldots, m\}\) if \(|\Omega| = m < \infty\) and \(\mathbb{Z}^+\) if \(\Omega\) is infinite, and let \(\sigma : \Omega^{\prime} \to \Omega\) be a bijection with \(s_{\sigma(1)} \ge s_{\sigma(2)} \ge \cdots\); replacing \(\Omega\), \(e_k\), \(h_k\), \(s_k\) by \(\Omega^{\prime}\), \(e_{\sigma(k)}\), \(h_{\sigma(k)}\), \(s_{\sigma(k)}\) changes neither the orthonormality of the two families nor the value of the unordered sum. Comparing with the definition 10.116 of the singular values, we may and do assume that \(\Omega\) is an initial segment of \(\mathbb{Z}^+\), that is, \(\Omega = \{1, \ldots, m\}\) for some \(m \in \mathbb{Z}^+\) or \(\Omega = \mathbb{Z}^+\) or \(\Omega = \emptyset\), and that

\begin{equation*} s_k = s_k(T) \text{ for } k \in \Omega, \qquad s_j(T) = 0 \text{ for } j \in \mathbb{Z}^+ \setminus \Omega . \end{equation*}

Taking \(f = e_j\) in the displayed formula for \(T\) gives

\begin{equation*} Te_j = s_j h_j \qquad \text{for each } j \in \Omega, \end{equation*}

because \(\langle e_j, e_k \rangle = 0\) for \(k \ne j\) and \(\langle e_j, e_j \rangle = 1\).

Let \(I\) denote the infimum in question.

Proof that \(I \le s_n(T)\). Define \(S \in \mathcal{B}(V)\) by

\begin{equation*} Sf = \sum_{k \in \Omega,\ k \le n - 1} s_k \langle f, e_k \rangle h_k . \end{equation*}

This is a finite sum of at most \(n - 1\) terms, so \(S\) is a bounded operator with \(\dim \operatorname{range} S \le n - 1 < n\). For \(f \in V\) we have

\begin{equation*} (T - S)f = \sum_{k \in \Omega,\ k \ge n} s_k \langle f, e_k \rangle h_k , \end{equation*}

so by 8.54(b) and then Bessel’s inequality (8.57),

\begin{equation*} \|(T - S)f\|^2 = \sum_{k \in \Omega,\ k \ge n} s_k^{\,2} |\langle f, e_k \rangle|^2 \le s_n(T)^2 \sum_{k \in \Omega} |\langle f, e_k \rangle|^2 \le s_n(T)^2 \|f\|^2 . \end{equation*}

Here the middle inequality is justified as follows. If \(n \in \Omega\), then \(s_k \le s_n = s_n(T)\) for every \(k \in \Omega\) with \(k \ge n\), because the \(s_k\) are in decreasing order. If \(n \notin \Omega\), then, since \(\Omega\) is an initial segment of \(\mathbb{Z}^+\), there is no \(k \in \Omega\) with \(k \ge n\); the sum on the left is then empty and both sides are \(0\). In either case we conclude that \(\|T - S\| \le s_n(T)\), and hence \(I \le s_n(T)\).

Proof that \(I \ge s_n(T)\). If \(s_n(T) = 0\) this is obvious, so assume \(s_n(T) > 0\); thus \(n \in \Omega\) and the orthonormal vectors \(e_1, \ldots, e_n\) are defined. Suppose \(S \in \mathcal{B}(V)\) with \(\dim \operatorname{range} S < n\). Let

\begin{equation*} U = \operatorname{span}\{e_1, \ldots, e_n\}, \end{equation*}

which has dimension \(n\) because \(e_1, \ldots, e_n\) are orthonormal and hence linearly independent. The restriction \(S|_U\) is a linear map from the \(n\)-dimensional space \(U\) into \(\operatorname{range} S\), whose dimension is at most \(n - 1\). By the Fundamental Theorem of Linear Maps (the rank-nullity theorem),

\begin{equation*} \dim \operatorname{null}(S|_U) = n - \dim \operatorname{range}(S|_U) \ge n - (n-1) = 1 . \end{equation*}

Hence there exists \(f \in U\) with \(\|f\| = 1\) and \(Sf = 0\). Write \(f = a_1 e_1 + \cdots + a_n e_n\); then \(\langle f, e_k \rangle = a_k\) for \(k \le n\) and \(\langle f, e_k \rangle = 0\) for \(k \in \Omega\) with \(k > n\), and \(\sum_{k=1}^{n} |a_k|^2 = \|f\|^2 = 1\). Therefore

\begin{equation*} Tf = \sum_{k=1}^{n} s_k a_k h_k , \end{equation*}

so, since \(\{h_k\}\) is orthonormal and \(s_k \ge s_n = s_n(T)\) for \(k \le n\),

\begin{equation*} \|(T - S)f\|^2 = \|Tf\|^2 = \sum_{k=1}^{n} s_k^{\,2} |a_k|^2 \ge s_n(T)^2 \sum_{k=1}^{n} |a_k|^2 = s_n(T)^2 . \end{equation*}

Because \(\|f\| = 1\), this gives \(\|T - S\| \ge s_n(T)\). As \(S\) was an arbitrary bounded operator with \(\dim \operatorname{range} S < n\), we conclude \(I \ge s_n(T)\).

Combining the two inequalities gives \(I = s_n(T)\), as desired. (Note that the first part of the proof shows the infimum is attained, by the truncated operator \(S\) constructed there.)

Problem (10D.16)

Suppose \(T\) is a compact operator on a Hilbert space \(V\) and \(n \in \mathbb{Z}^+\). Prove that

\begin{equation*} s_n(T) = \inf\{\|T|_{U^{\perp}}\| : U \text{ is a subspace of } V \text{ with } \dim U < n\}. \end{equation*}

Solution

As in the solution to Exercise 10D.15, fix a singular value decomposition of \(T\) (10.113),

\begin{equation*} Tf = \sum_{k \in \Omega} s_k \langle f, e_k \rangle h_k \qquad \text{for all } f \in V, \end{equation*}

where \(\{e_k\}_{k \in \Omega}\) and \(\{h_k\}_{k \in \Omega}\) are orthonormal families, and where (using the uniqueness and decreasing rearrangement of the singular values discussed before 10.116) \(\Omega\) is an initial segment of \(\mathbb{Z}^+\), \(s_k = s_k(T)\) for \(k \in \Omega\), and \(s_j(T) = 0\) for \(j \in \mathbb{Z}^+ \setminus \Omega\). Let \(I\) denote the infimum on the right side of the equation to be proved; here \(\|T|_{U^{\perp}}\| = \sup\{\|Tf\| : f \in U^{\perp} \text{ and } \|f\| \le 1\}\).

Proof that \(I \le s_n(T)\). Let

\begin{equation*} U = \operatorname{span}\{e_k : k \in \Omega \text{ and } k \le n - 1\}, \end{equation*}

so \(\dim U \le n - 1 < n\). If \(f \in U^{\perp}\), then \(\langle f, e_k \rangle = 0\) for every \(k \in \Omega\) with \(k \le n-1\), and hence

\begin{equation*} Tf = \sum_{k \in \Omega,\ k \ge n} s_k \langle f, e_k \rangle h_k . \end{equation*}

Because \(\{h_k\}_{k \in \Omega}\) is orthonormal, 8.54(b) gives

\begin{equation*} \|Tf\|^2 = \sum_{k \in \Omega,\ k \ge n} s_k^{\,2} |\langle f, e_k \rangle|^2 \le s_n(T)^2 \sum_{k \in \Omega} |\langle f, e_k \rangle|^2 \le s_n(T)^2 \|f\|^2 , \end{equation*}

where the last step is Bessel’s inequality (8.57) and the middle step holds because \(s_k \le s_n = s_n(T)\) for every \(k \in \Omega\) with \(k \ge n\) when \(n \in \Omega\), while if \(n \notin \Omega\) there is no such \(k\) (as \(\Omega\) is an initial segment) and the sum is empty. Thus \(\|T|_{U^{\perp}}\| \le s_n(T)\), and therefore \(I \le s_n(T)\).

Proof that \(I \ge s_n(T)\). If \(s_n(T) = 0\) there is nothing to prove, so assume \(s_n(T) > 0\); then \(n \in \Omega\), so \(e_1, \ldots, e_n\) are defined and orthonormal. Let \(U\) be any subspace of \(V\) with \(d := \dim U < n\), and let \(u_1, \ldots, u_d\) be a basis of \(U\). Let

\begin{equation*} W = \operatorname{span}\{e_1, \ldots, e_n\}, \end{equation*}

a subspace of dimension \(n\). Define a linear map \(\varphi : W \to \mathbb{F}^d\) by

\begin{equation*} \varphi(f) = (\langle f, u_1 \rangle, \ldots, \langle f, u_d \rangle). \end{equation*}

Because \(\dim W = n > d \ge \dim \operatorname{range} \varphi\), the Fundamental Theorem of Linear Maps implies \(\operatorname{null} \varphi \ne \{0\}\). Choose \(f \in \operatorname{null} \varphi\) with \(\|f\| = 1\). Then \(f\) is orthogonal to each \(u_j\), hence (by linearity of \(g \mapsto \langle f, g \rangle\) in the appropriate slot, or simply because every element of \(U\) is a linear combination of \(u_1, \ldots, u_d\)) \(f\) is orthogonal to every element of \(U\); that is, \(f \in U^{\perp}\).

Write \(f = a_1 e_1 + \cdots + a_n e_n\), so \(\sum_{k=1}^{n} |a_k|^2 = \|f\|^2 = 1\), and \(\langle f, e_k \rangle = a_k\) for \(k \le n\) while \(\langle f, e_k \rangle = 0\) for \(k \in \Omega\) with \(k > n\). Hence

\begin{equation*} \|Tf\|^2 = \Bigl\| \sum_{k=1}^{n} s_k a_k h_k \Bigr\|^2 = \sum_{k=1}^{n} s_k^{\,2}|a_k|^2 \ge s_n(T)^2 \sum_{k=1}^{n} |a_k|^2 = s_n(T)^2 , \end{equation*}

using \(s_k \ge s_n = s_n(T)\) for \(k \le n\). Since \(f \in U^{\perp}\) and \(\|f\| = 1\), this shows \(\|T|_{U^{\perp}}\| \ge s_n(T)\). As \(U\) was an arbitrary subspace of dimension less than \(n\), we get \(I \ge s_n(T)\).

Combining the two inequalities gives \(s_n(T) = I\), as desired. (The first part of the proof shows that the infimum is attained.)

Problem (10D.17)

Suppose \(T\) is a compact operator on a Hilbert space \(V\) with singular value decomposition

\begin{equation*} Tf = \sum_{k \in \Omega} s_k \langle f, e_k \rangle h_k \end{equation*}

for all \(f \in V\). Prove that

\begin{equation*} T^* f = \sum_{k \in \Omega} s_k \langle f, h_k \rangle e_k \end{equation*}

for all \(f \in V\).

Solution

Here \(\Omega\) is a countable set, \(\{e_k\}_{k \in \Omega}\) and \(\{h_k\}_{k \in \Omega}\) are orthonormal families in \(V\), and \(\{s_k\}_{k \in \Omega}\) is a family of positive numbers, as in 10.113.

Step 1: the numbers \(s_k\) are uniformly bounded. Taking \(f = e_j\) in the formula for \(T\) and using the orthonormality of \(\{e_k\}_{k \in \Omega}\) gives \(Te_j = s_j h_j\) for each \(j \in \Omega\). Because \(\|h_j\| = \|e_j\| = 1\), this implies

\begin{equation*} s_j = \|s_j h_j\| = \|Te_j\| \le \|T\| \qquad \text{for every } j \in \Omega . \end{equation*}

Step 2: the asserted formula defines a bounded operator. For \(f \in V\), Bessel’s inequality (8.57) applied to the orthonormal family \(\{h_k\}_{k \in \Omega}\) gives

\begin{equation*} \sum_{k \in \Omega} |s_k \langle f, h_k \rangle|^2 \le \|T\|^2 \sum_{k \in \Omega} |\langle f, h_k \rangle|^2 \le \|T\|^2 \|f\|^2 < \infty . \end{equation*}

Hence by 8.54(a) the unordered sum

\begin{equation*} Sf := \sum_{k \in \Omega} s_k \langle f, h_k \rangle e_k \end{equation*}

converges in \(V\), and by 8.54(b) we have \(\|Sf\|^2 = \sum_{k \in \Omega}|s_k \langle f, h_k \rangle|^2 \le \|T\|^2 \|f\|^2\). The map \(S : V \to V\) so defined is linear (the coefficients \(f \mapsto s_k \langle f, h_k \rangle\) depend linearly on \(f\), and unordered sums are additive and homogeneous), and the estimate just obtained shows that \(S\) is bounded, with \(\|S\| \le \|T\|\).

Step 3: a remark on passing a continuous linear functional through an unordered sum. Suppose \(\sum_{k \in \Omega} x_k = x\) in \(V\) (unordered convergence) and \(\varphi\) is a bounded linear functional on \(V\). Let \(\varepsilon > 0\). There is a finite set \(F_0 \subseteq \Omega\) such that \(\|x - \sum_{k \in F} x_k\| < \varepsilon / (\|\varphi\| + 1)\) for every finite \(F\) with \(F_0 \subseteq F \subseteq \Omega\); for such \(F\),

\begin{equation*} \Bigl| \varphi(x) - \sum_{k \in F} \varphi(x_k) \Bigr| = \Bigl| \varphi\Bigl(x - \sum_{k \in F} x_k\Bigr) \Bigr| < \varepsilon . \end{equation*}

Thus \(\sum_{k \in \Omega} \varphi(x_k) = \varphi(x)\).

Step 4: \(S\) satisfies the defining property of \(T^*\). Let \(f, g \in V\). The functional \(x \mapsto \langle x, g \rangle\) is linear and bounded (Cauchy–Schwarz), so Step 3 applied to the convergent sum defining \(Tf\) gives

\begin{equation*} \langle Tf, g \rangle = \sum_{k \in \Omega} s_k \langle f, e_k \rangle \overline{\langle g, h_k \rangle} . \end{equation*}

Similarly, applying Step 3 to the convergent sum defining \(Sg\) with the functional \(x \mapsto \langle x, f \rangle\) gives

\begin{equation*} \langle Sg, f \rangle = \sum_{k \in \Omega} s_k \langle g, h_k \rangle \overline{\langle f, e_k \rangle}, \end{equation*}

and hence, taking complex conjugates (which is legitimate term by term for a convergent unordered sum of scalars),

\begin{equation*} \langle f, Sg \rangle = \overline{\langle Sg, f \rangle} = \sum_{k \in \Omega} s_k \overline{\langle g, h_k \rangle} \langle f, e_k \rangle . \end{equation*}

The two displayed sums are identical. Therefore

\begin{equation*} \langle Tf, g \rangle = \langle f, Sg \rangle \qquad \text{for all } f, g \in V . \end{equation*}

Step 5: conclusion. By the definition of the adjoint (10.1) we also have \(\langle Tf, g \rangle = \langle f, T^*g \rangle\) for all \(f, g \in V\). Subtracting, \(\langle f, T^*g - Sg \rangle = 0\) for all \(f \in V\) and all \(g \in V\). Taking \(f = T^*g - Sg\) gives \(\|T^*g - Sg\|^2 = 0\), so \(T^*g = Sg\) for every \(g \in V\). That is,

\begin{equation*} T^* f = \sum_{k \in \Omega} s_k \langle f, h_k \rangle e_k \end{equation*}

for all \(f \in V\), as desired. (In particular this exhibits a singular value decomposition of \(T^*\) with the roles of the two orthonormal families interchanged, so \(T\) and \(T^*\) have the same singular values.)

Problem (10D.18)

Suppose that \(T\) is an operator on a finite-dimensional Hilbert space \(V\) with \(\dim V = n\).

(a) Prove that \(T\) is invertible if and only if \(s_n(T) \ne 0\).

(b) Suppose \(T\) is invertible and \(T\) has a singular value decomposition

\begin{equation*} Tf = s_1(T) \langle f, e_1 \rangle h_1 + \cdots + s_n(T) \langle f, e_n \rangle h_n \end{equation*}

for all \(f \in V\). Show that

\begin{equation*} T^{-1} f = \frac{\langle f, h_1 \rangle}{s_1(T)} e_1 + \cdots + \frac{\langle f, h_n \rangle}{s_n(T)} e_n \end{equation*}

for all \(f \in V\).

Solution

Every linear map on a finite-dimensional normed vector space is bounded, and every bounded operator with finite-dimensional range is compact (a bounded subset of a finite-dimensional space has compact closure). Thus \(T\) is a compact operator and the theory of singular values applies to it.

(a) By 10.113 and the decreasing rearrangement of the singular values discussed before 10.116, we may write

\begin{equation*} Tf = \sum_{k \in \Omega} s_k(T) \langle f, e_k \rangle h_k \qquad \text{for all } f \in V, \end{equation*}

where \(\{e_k\}_{k \in \Omega}\) and \(\{h_k\}_{k \in \Omega}\) are orthonormal families in \(V\), where \(\Omega\) is an initial segment of \(\mathbb{Z}^+\) (possibly empty), where \(s_k(T) > 0\) for \(k \in \Omega\), and where \(s_j(T) = 0\) for \(j \in \mathbb{Z}^+ \setminus \Omega\).

Because \(\{e_k\}_{k \in \Omega}\) is an orthonormal, hence linearly independent, family in \(V\), we have \(|\Omega| \le \dim V = n\); in particular \(\Omega\) is finite.

The formula above shows that \(\operatorname{range} T \subseteq \operatorname{span}\{h_k\}_{k \in \Omega}\). Conversely, taking \(f = e_j\) gives \(Te_j = s_j(T) h_j\), so \(h_j = Te_j / s_j(T) \in \operatorname{range} T\) for each \(j \in \Omega\). Hence

\begin{equation*} \operatorname{range} T = \operatorname{span}\{h_k\}_{k \in \Omega}, \qquad \dim \operatorname{range} T = |\Omega|, \end{equation*}

the last equality because \(\{h_k\}_{k \in \Omega}\) is orthonormal and therefore linearly independent.

Now, since \(\Omega\) is an initial segment of \(\mathbb{Z}^+\) with \(|\Omega| \le n\), we have

\begin{equation*} s_n(T) \ne 0 \iff n \in \Omega \iff |\Omega| = n \iff \dim \operatorname{range} T = n = \dim V \iff T \text{ is surjective}. \end{equation*}

Finally, an operator on a finite-dimensional vector space is invertible if and only if it is surjective (surjectivity forces injectivity by the Fundamental Theorem of Linear Maps, and the inverse of a linear bijection of a finite-dimensional normed space is automatically linear and bounded). Therefore \(T\) is invertible if and only if \(s_n(T) \ne 0\).

(b) Suppose \(T\) is invertible with the singular value decomposition displayed in the statement. By part (a), \(s_n(T) \ne 0\), and since \(s_1(T) \ge \cdots \ge s_n(T) > 0\), all of \(s_1(T), \ldots, s_n(T)\) are positive; in particular the formula defining \(T^{-1}\) below makes sense.

The vectors \(e_1, \ldots, e_n\) are orthonormal, hence linearly independent, and there are \(n = \dim V\) of them; therefore \(\{e_1, \ldots, e_n\}\) is a basis of \(V\). The same reasoning shows that \(\{h_1, \ldots, h_n\}\) is a basis of \(V\).

Define \(R : V \to V\) by

\begin{equation*} Rf = \frac{\langle f, h_1 \rangle}{s_1(T)} e_1 + \cdots + \frac{\langle f, h_n \rangle}{s_n(T)} e_n , \end{equation*}

which is a linear map. Because \(\{e_k\}\) and \(\{h_k\}\) are orthonormal, evaluating the two formulas on basis vectors gives

\begin{equation*} Te_j = s_j(T) h_j \quad \text{and} \quad Rh_j = \frac{e_j}{s_j(T)} \qquad \text{for } j = 1, \ldots, n . \end{equation*}

Hence for each \(j\),

\begin{equation*} (RT)e_j = R\bigl(s_j(T) h_j\bigr) = s_j(T) \cdot \frac{e_j}{s_j(T)} = e_j, \qquad (TR)h_j = T\Bigl(\frac{e_j}{s_j(T)}\Bigr) = \frac{s_j(T) h_j}{s_j(T)} = h_j . \end{equation*}

Since \(RT\) and the identity agree on the basis \(\{e_1, \ldots, e_n\}\) and both are linear, \(RT = I\); since \(TR\) and the identity agree on the basis \(\{h_1, \ldots, h_n\}\), \(TR = I\). Therefore \(R = T^{-1}\), which is exactly the asserted formula

\begin{equation*} T^{-1} f = \frac{\langle f, h_1 \rangle}{s_1(T)} e_1 + \cdots + \frac{\langle f, h_n \rangle}{s_n(T)} e_n \end{equation*}

for all \(f \in V\). (Note that this formula is a singular value decomposition of \(T^{-1}\) once the terms are rewritten in decreasing order of the coefficients \(1/s_k(T)\); thus \(s_k(T^{-1}) = 1/s_{n+1-k}(T)\) for \(k = 1, \ldots, n\).)

Problem (10D.19)

Suppose \(T\) is a compact operator on a Hilbert space \(V\). Prove that

\begin{equation*} \sum_{k \in \Gamma} \|Te_k\|^2 = \sum_{n=1}^{\infty} \bigl(s_n(T)\bigr)^2 \end{equation*}

for every orthonormal basis \(\{e_k\}_{k \in \Gamma}\) of \(V\).

Solution

All sums below have nonnegative terms, so they are interpreted as elements of \([0, \infty]\): for a family \(\{a_i\}_{i \in I}\) in \([0, \infty)\) we set

\begin{equation*} \sum_{i \in I} a_i = \sup\Bigl\{ \sum_{i \in F} a_i : F \subseteq I \text{ and } F \text{ is finite} \Bigr\}, \end{equation*}

which agrees with the unordered sum when that sum converges. The asserted equality is an equality in \([0, \infty]\) (both sides are \(+\infty\) for some compact operators).

Step 1: interchange of nonnegative double sums. Suppose \(\{a_{k,j}\}_{k \in \Gamma, j \in \Lambda}\) is a family in \([0, \infty)\). Let \(A = \sum_{k \in \Gamma} \sum_{j \in \Lambda} a_{k,j}\) and \(B = \sum_{j \in \Lambda} \sum_{k \in \Gamma} a_{k,j}\). If \(F \subseteq \Gamma\) and \(G \subseteq \Lambda\) are finite, then

\begin{equation*} \sum_{j \in G} \sum_{k \in F} a_{k,j} = \sum_{k \in F} \sum_{j \in G} a_{k,j} \le \sum_{k \in F} \sum_{j \in \Lambda} a_{k,j} \le A , \end{equation*}

the first equality because both sides are finite sums of the same finitely many nonnegative numbers.

Now fix a finite \(G \subseteq \Lambda\) and let \(c\) be any real number with \(c < \sum_{j \in G} \sum_{k \in \Gamma} a_{k,j}\). Since \(G\) is finite, we may choose real numbers \(t_j\) for \(j \in G\) with \(t_j < \sum_{k \in \Gamma} a_{k,j}\) for each \(j\) and \(\sum_{j \in G} t_j > c\). By the definition of the inner sums as suprema, for each \(j \in G\) there is a finite \(F_j \subseteq \Gamma\) with \(\sum_{k \in F_j} a_{k,j} > t_j\). Put \(F = \bigcup_{j \in G} F_j\), a finite subset of \(\Gamma\). Because all terms are nonnegative, \(\sum_{k \in F} a_{k,j} \ge \sum_{k \in F_j} a_{k,j} > t_j\) for each \(j \in G\), so

\begin{equation*} \sum_{j \in G} \sum_{k \in F} a_{k,j} > \sum_{j \in G} t_j > c . \end{equation*}

Combined with the previous display, \(c < A\). As \(c\) was any real number below \(\sum_{j \in G} \sum_{k \in \Gamma} a_{k,j}\), we get \(\sum_{j \in G} \sum_{k \in \Gamma} a_{k,j} \le A\). Taking the supremum over finite \(G \subseteq \Lambda\) gives \(B \le A\). By symmetry \(A \le B\), so \(A = B\).

Step 2: the quantity \(\sum_{k \in \Gamma} \|Te_k\|^2\) is the same for every orthonormal basis. Let \(\{e_k\}_{k \in \Gamma}\) and \(\{g_j\}_{j \in \Lambda}\) be orthonormal bases of \(V\). By Parseval’s identity [8.63(c)] applied to the basis \(\{g_j\}_{j \in \Lambda}\) and to the basis \(\{e_k\}_{k \in \Gamma}\), together with Step 1 and the definition of the adjoint (10.1),

\begin{equation*} \begin{aligned} \sum_{k \in \Gamma} \|Te_k\|^2 &= \sum_{k \in \Gamma} \sum_{j \in \Lambda} |\langle Te_k, g_j \rangle|^2 \\ &= \sum_{j \in \Lambda} \sum_{k \in \Gamma} |\langle e_k, T^* g_j \rangle|^2 \\ &= \sum_{j \in \Lambda} \|T^* g_j\|^2 . \end{aligned} \end{equation*}

The last expression does not involve \(\{e_k\}_{k \in \Gamma}\). Hence \(\sum_{k \in \Gamma} \|Te_k\|^2\) has the same value for every orthonormal basis \(\{e_k\}_{k \in \Gamma}\) of \(V\).

Step 3: computing the common value from a singular value decomposition. By 10.113 and the decreasing rearrangement of singular values discussed before 10.116, write

\begin{equation*} Tf = \sum_{k \in \Omega} s_k(T) \langle f, u_k \rangle v_k \qquad \text{for all } f \in V, \end{equation*}

where \(\{u_k\}_{k \in \Omega}\) and \(\{v_k\}_{k \in \Omega}\) are orthonormal families in \(V\), \(\Omega\) is an initial segment of \(\mathbb{Z}^+\), \(s_k(T) > 0\) for \(k \in \Omega\), and \(s_j(T) = 0\) for \(j \in \mathbb{Z}^+ \setminus \Omega\).

Let \(U = \overline{\operatorname{span}}\{u_k\}_{k \in \Omega}\), a closed subspace of \(V\); thus \(\{u_k\}_{k \in \Omega}\) is an orthonormal basis of the Hilbert space \(U\) (8.61). The subspace \(U^{\perp}\) is closed, hence is itself a Hilbert space, so by 8.75 it has an orthonormal basis \(\{w_i\}_{i \in \Lambda}\). The family \(\{u_k\}_{k \in \Omega} \cup \{w_i\}_{i \in \Lambda}\) is orthonormal (each \(w_i\) is orthogonal to each \(u_k\) because \(u_k \in U\)), and the closure of its span contains \(U\) and contains \(U^{\perp}\), hence contains \(U + U^{\perp} = V\) by the orthogonal decomposition 8.43. Therefore this family is an orthonormal basis of \(V\).

Evaluate \(T\) on it. Taking \(f = u_j\) in the formula for \(T\) gives \(Tu_j = s_j(T) v_j\), so

\begin{equation*} \|Tu_j\|^2 = s_j(T)^2 \qquad \text{for } j \in \Omega , \end{equation*}

because \(\|v_j\| = 1\). If \(i \in \Lambda\), then \(w_i \in U^{\perp}\), so \(\langle w_i, u_k \rangle = 0\) for every \(k \in \Omega\) and hence \(Tw_i = 0\).

Consequently, since the sum of a family of nonnegative numbers indexed by a disjoint union is the sum of the two sub-sums (immediate from the supremum definition, because every finite subset of \(\Omega \sqcup \Lambda\) splits into a finite subset of \(\Omega\) and a finite subset of \(\Lambda\)), the sum in question for this particular orthonormal basis is

\begin{equation*} \sum_{k \in \Omega} \|Tu_k\|^2 + \sum_{i \in \Lambda} \|Tw_i\|^2 = \sum_{k \in \Omega} s_k(T)^2 + 0 = \sum_{n=1}^{\infty} \bigl(s_n(T)\bigr)^2 , \end{equation*}

where the last equality holds because \(s_n(T) = 0\) for every \(n \in \mathbb{Z}^+ \setminus \Omega\), so the omitted terms contribute nothing (and for a family of nonnegative numbers indexed by a subset of \(\mathbb{Z}^+\), the unordered sum agrees with the sum of the series, since all terms are nonnegative).

By Step 2, the value just computed equals \(\sum_{k \in \Gamma} \|Te_k\|^2\) for every orthonormal basis \(\{e_k\}_{k \in \Gamma}\) of \(V\). Therefore

\begin{equation*} \sum_{k \in \Gamma} \|Te_k\|^2 = \sum_{n=1}^{\infty} \bigl(s_n(T)\bigr)^2 \end{equation*}

for every orthonormal basis \(\{e_k\}_{k \in \Gamma}\) of \(V\), as desired. (Taking \(T = \mathcal{I}_K\) and using this with the standard computation recovers 10.120.)

Problem (10D.20)

Use the result of Example 10.124 to evaluate \(\displaystyle\sum_{n=1}^{\infty} \frac{1}{n^2}\).

Solution

Example 10.124 states that

\begin{equation*} \frac{1}{1^2} + \frac{1}{3^2} + \frac{1}{5^2} + \cdots = \frac{\pi^2}{8}, \end{equation*}

in other words \(\sum_{k=1}^{\infty} \frac{1}{(2k-1)^2} = \frac{\pi^2}{8}\).

First note that \(\sum_{n=1}^{\infty} \frac{1}{n^2}\) converges: its terms are positive, and for every \(N \in \mathbb{Z}^+\),

\begin{equation*} \sum_{n=1}^{N} \frac{1}{n^2} \le 1 + \sum_{n=2}^{N} \frac{1}{n(n-1)} = 1 + \sum_{n=2}^{N} \Bigl( \frac{1}{n-1} - \frac{1}{n} \Bigr) = 2 - \frac{1}{N} < 2 , \end{equation*}

so the increasing sequence of partial sums is bounded above and hence converges. Let

\begin{equation*} S = \sum_{n=1}^{\infty} \frac{1}{n^2} \in (0, \infty). \end{equation*}

Now split the partial sum over the first \(2N\) integers into its odd-indexed and even-indexed terms:

\begin{equation*} \sum_{n=1}^{2N} \frac{1}{n^2} = \sum_{k=1}^{N} \frac{1}{(2k-1)^2} + \sum_{k=1}^{N} \frac{1}{(2k)^2} = \sum_{k=1}^{N} \frac{1}{(2k-1)^2} + \frac{1}{4} \sum_{k=1}^{N} \frac{1}{k^2} . \end{equation*}

Let \(N \to \infty\). The left side converges to \(S\) (it is a subsequence of the convergent sequence of partial sums), the first sum on the right converges to \(\pi^2/8\) by Example 10.124, and the second converges to \(S/4\). Hence

\begin{equation*} S = \frac{\pi^2}{8} + \frac{S}{4}, \end{equation*}

which gives \(\frac{3}{4} S = \frac{\pi^2}{8}\) and therefore

\begin{equation*} \sum_{n=1}^{\infty} \frac{1}{n^2} = \frac{4}{3} \cdot \frac{\pi^2}{8} = \frac{\pi^2}{6} . \end{equation*}

Problem (10D.21)

Suppose \(T\) is a normal compact operator on a complex Hilbert space. Prove that the following are equivalent.

(a) \(\operatorname{range} T\) is finite-dimensional.

(b) \(\operatorname{sp}(T)\) is a finite set.

(c) \(s_n(T) = 0\) for some \(n \in \mathbb{Z}\).

Solution

Let \(V\) denote the complex Hilbert space on which \(T\) acts. The book prints \(n \in \mathbb{Z}\) in (c); since 10.116 defines \(s_n(T)\) only for \(n \in \mathbb{Z}^+\), condition (c) is to be read as: \(s_n(T) = 0\) for some \(n \in \mathbb{Z}^+\). We prove (a) \(\iff\) (c), then (a) \(\implies\) (b), then (b) \(\implies\) (a).

(a) \(\iff\) (c). This part uses only compactness. By 10.113 and the decreasing rearrangement of the singular values discussed before 10.116, write

\begin{equation*} Tf = \sum_{k \in \Omega} s_k(T) \langle f, u_k \rangle v_k \qquad \text{for all } f \in V, \end{equation*}

where \(\{u_k\}_{k \in \Omega}\) and \(\{v_k\}_{k \in \Omega}\) are orthonormal families in \(V\), \(\Omega\) is an initial segment of \(\mathbb{Z}^+\) (that is, \(\Omega = \emptyset\), or \(\Omega = \{1, \ldots, m\}\), or \(\Omega = \mathbb{Z}^+\)), \(s_k(T) > 0\) for \(k \in \Omega\), and \(s_j(T) = 0\) for \(j \in \mathbb{Z}^+ \setminus \Omega\).

Taking \(f = u_j\) gives \(Tu_j = s_j(T) v_j\), so \(v_j = Tu_j / s_j(T) \in \operatorname{range} T\) for each \(j \in \Omega\).

If \(\Omega\) is infinite, then \(\operatorname{range} T\) contains the infinite orthonormal family \(\{v_k\}_{k \in \Omega}\), which is linearly independent; hence \(\operatorname{range} T\) is infinite-dimensional. If \(\Omega\) is finite, then the displayed formula shows \(\operatorname{range} T \subseteq \operatorname{span}\{v_k\}_{k \in \Omega}\), so \(\operatorname{range} T\) is finite-dimensional (indeed \(\dim \operatorname{range} T = |\Omega|\), by the previous paragraph). Thus

\begin{equation*} \operatorname{range} T \text{ is finite-dimensional} \iff \Omega \text{ is finite}. \end{equation*}

Finally, \(\Omega\) is finite if and only if \(s_n(T) = 0\) for some \(n \in \mathbb{Z}^+\): if \(\Omega\) is finite take any \(n > \max \Omega\) (or \(n = 1\) if \(\Omega = \emptyset\)), and conversely if \(s_n(T) = 0\) then \(n \notin \Omega\), which forces \(\Omega \subseteq \{1, \ldots, n-1\}\) because \(\Omega\) is an initial segment of \(\mathbb{Z}^+\). This proves (a) \(\iff\) (c).

(a) \(\implies\) (b). Suppose \(m := \dim \operatorname{range} T < \infty\). Suppose \(\alpha_1, \ldots, \alpha_j\) are distinct nonzero eigenvalues of \(T\), with corresponding eigenvectors \(f_1, \ldots, f_j\) (each nonzero). Then

\begin{equation*} f_i = \frac{T f_i}{\alpha_i} \in \operatorname{range} T \qquad \text{for } i = 1, \ldots, j . \end{equation*}

Because \(T\) is normal, 10.57 implies that \(f_1, \ldots, f_j\) are pairwise orthogonal; being nonzero and pairwise orthogonal, they are linearly independent. Since they all lie in the \(m\)-dimensional subspace \(\operatorname{range} T\), we get \(j \le m\). Hence \(T\) has at most \(m\) distinct nonzero eigenvalues.

Now let \(\alpha \in \operatorname{sp}(T)\) with \(\alpha \ne 0\). The Fredholm Alternative (10.85) implies that \(\alpha\) is an eigenvalue of \(T\). Therefore

\begin{equation*} \operatorname{sp}(T) \subseteq \{\text{nonzero eigenvalues of } T\} \cup \{0\}, \end{equation*}

and the right side has at most \(m + 1\) elements. Thus \(\operatorname{sp}(T)\) is a finite set.

(b) \(\implies\) (a). Suppose \(\operatorname{sp}(T)\) is a finite set. Because \(V\) is a complex Hilbert space and \(T\) is a normal compact operator, the Spectral Theorem for normal compact operators (10.107) provides an orthonormal basis \(\{f_k\}_{k \in \Gamma}\) of \(V\) consisting of eigenvectors of \(T\); say

\begin{equation*} T f_k = \lambda_k f_k \qquad \text{for each } k \in \Gamma , \end{equation*}

with \(\lambda_k \in \mathbb{C}\).

Each \(\lambda_k\) is an eigenvalue of \(T\), so \(T - \lambda_k I\) is not injective and hence not invertible; thus \(\lambda_k \in \operatorname{sp}(T)\). Consequently the set of values \(\{\lambda_k : k \in \Gamma\}\) is contained in the finite set \(\operatorname{sp}(T)\), so only finitely many distinct values occur among the \(\lambda_k\).

Fix a value \(\lambda \ne 0\) occurring among the \(\lambda_k\). The family \(\{f_k : \lambda_k = \lambda\}\) is an orthonormal family contained in \(\operatorname{null}(T - \lambda I)\), which is finite-dimensional by 10.82 (nonzero eigenvalues of compact operators have finite multiplicity). An orthonormal family is linearly independent, so \(\{k \in \Gamma : \lambda_k = \lambda\}\) is finite. Because only finitely many distinct nonzero values \(\lambda\) occur, the set

\begin{equation*} \Lambda = \{k \in \Gamma : \lambda_k \ne 0\} \end{equation*}

is a finite union of finite sets and hence is finite.

Now let \(f \in V\). By Parseval’s identity [8.63(a)], \(f = \sum_{k \in \Gamma} \langle f, f_k \rangle f_k\), with the unordered sum converging in \(V\). Since \(T\) is a bounded linear operator, it may be applied term by term: if \(\sum_{k \in \Gamma} x_k = x\) and \(\varepsilon > 0\), choose a finite \(F_0 \subseteq \Gamma\) with \(\|x - \sum_{k \in F} x_k\| < \varepsilon/(\|T\| + 1)\) for all finite \(F\) with \(F_0 \subseteq F \subseteq \Gamma\); then \(\|Tx - \sum_{k \in F} Tx_k\| < \varepsilon\) for all such \(F\), so \(\sum_{k \in \Gamma} Tx_k = Tx\). Applying this,

\begin{equation*} Tf = \sum_{k \in \Gamma} \langle f, f_k \rangle T f_k = \sum_{k \in \Gamma} \lambda_k \langle f, f_k \rangle f_k = \sum_{k \in \Lambda} \lambda_k \langle f, f_k \rangle f_k , \end{equation*}

the last equality because the terms with \(k \notin \Lambda\) are \(0\). Hence

\begin{equation*} \operatorname{range} T \subseteq \operatorname{span}\{f_k\}_{k \in \Lambda}, \end{equation*}

which is finite-dimensional because \(\Lambda\) is finite. Thus \(\operatorname{range} T\) is finite-dimensional, proving (a).

Having shown (a) \(\iff\) (c), (a) \(\implies\) (b), and (b) \(\implies\) (a), we conclude that (a), (b), and (c) are equivalent.

Note that normality is needed only in the implication (b) \(\implies\) (a): for a general compact operator, (a) and (c) are still equivalent, and (a) still implies (b) (linear independence of eigenvectors corresponding to distinct eigenvalues does not require normality), but the Volterra operator \(V\) shows that (b) does not imply (a) in general: by 10.89 we have \(\operatorname{sp}(V) = \{0\}\), a finite set, while \(V\) is injective on the infinite-dimensional space \(L^2([0,1])\) and hence has infinite-dimensional range.

Problem (10D.22)

Find the singular values of the Volterra operator.

[Your answer, when combined with Exercise 12, should show that the norm of the Volterra operator is \(\frac{2}{\pi}\). This appearance of \(\pi\) can be surprising because the definition of the Volterra operator does not involve \(\pi\).]

Solution

Recall (Example 10.15) that the Volterra operator \(V : L^2([0,1]) \to L^2([0,1])\) is defined by

\begin{equation*} (Vf)(x) = \int_0^x f . \end{equation*}

We will show that the singular values of \(V\) are

\begin{equation*} s_n(V) = \frac{2}{(2n-1)\pi} \qquad \text{for } n \in \mathbb{Z}^+, \end{equation*}

that is, the list of singular values is

\begin{equation*} \frac{2}{\pi} \ge \frac{2}{3\pi} \ge \frac{2}{5\pi} \ge \frac{2}{7\pi} \ge \cdots , \end{equation*}

each value occurring exactly once. Everything below works for \(\mathbb{F} = \mathbb{R}\) and for \(\mathbb{F} = \mathbb{C}\).

Step 1: a formula for \(V^\ast\).

The operator \(V\) is the integral operator \(\mathcal{I}_K\) associated with the kernel

\begin{equation*} K(x,y) = \begin{cases} 1 & \text{if } y < x, \\ 0 & \text{if } y \ge x, \end{cases} \end{equation*}

which lies in \(L^2(\mu \times \mu)\) for \(\mu\) Lebesgue measure on \([0,1]\); hence \(V\) is compact (see the paragraph following the proof of 10.70). By 10.16 the adjoint of \(\mathcal{I}_K\) is the integral operator with kernel \((x,y) \mapsto \overline{K(y,x)}\), and therefore

\begin{equation*} (V^\ast f)(x) = \int_x^1 f . \end{equation*}

(Directly: for \(f, g \in L^2([0,1])\), Tonelli/Fubini (5.28, 5.32) gives \(\langle Vf, g\rangle = \int_0^1 \int_0^1 \chi_{\{y<x\}} f(y) \overline{g(x)} \, dy \, dx = \int_0^1 f(y) \overline{\int_y^1 g} \, dy = \langle f, V^\ast g\rangle\) with \(V^\ast\) as displayed.)

Consequently

\begin{equation*} (V^\ast V f)(x) = \int_x^1 \Big( \int_0^t f \Big) \, dt . \end{equation*}

Step 2: regularity of an eigenvector of \(V^\ast V\).

Because \(V\) is compact, so is \(V^\ast\) (by 10.73), and hence \(V^\ast V\) is a compact self-adjoint operator. By the definition 10.116, the singular values of \(V\) are the positive square roots of the positive eigenvalues of \(V^\ast V\), arranged in decreasing order, with each singular value \(s\) listed as many times as the geometric multiplicity of \(s^2\) as an eigenvalue of \(V^\ast V\). So suppose \(\alpha > 0\) and \(f \in L^2([0,1])\) satisfies \(f \ne 0\) and

\begin{equation*} V^\ast V f = \alpha f . \end{equation*}

Since \([0,1]\) has finite measure, Hölder’s inequality (7.9) applied to \(|f| \cdot 1\) gives \(\|f\|_1 \le \|f\|_2\), so \(L^2([0,1]) \subset L^1([0,1])\) and \(Vf\) is well defined; moreover \(Vf\) is continuous on \([0,1]\), because \(|(Vf)(x) - (Vf)(x^{\prime})| \le \int_{[x^{\prime},x]} |f| \to 0\) as \(x^{\prime} \to x\) by the Dominated Convergence Theorem (3.31). Because \(Vf\) is continuous, the Fundamental Theorem of Calculus shows that \(V^\ast V f\) is continuously differentiable on \([0,1]\) with

\begin{equation*} (V^\ast V f)^{\prime}(x) = -(Vf)(x) \qquad \text{for all } x \in [0,1]. \end{equation*}

Since \(\alpha \ne 0\), the equation \(f = \frac{1}{\alpha} V^\ast V f\) shows that \(f\) agrees almost everywhere with a continuously differentiable function on \([0,1]\); replace \(f\) by that representative (this changes nothing in \(L^2([0,1])\)). Now \(f\) is continuous, so the Fundamental Theorem of Calculus gives \((Vf)^{\prime} = f\) everywhere on \([0,1]\). Hence \(V^\ast V f\) is twice continuously differentiable, and therefore so is \(f\), with

\begin{equation*} \alpha f^{\prime\prime}(x) = (V^\ast V f)^{\prime\prime}(x) = -(Vf)^{\prime}(x) = -f(x) \qquad \text{for all } x \in [0,1]. \end{equation*}

Step 3: the boundary conditions.

After the replacement made in Step 2, both sides of \(V^\ast V f = \alpha f\) are continuous functions on \([0,1]\) that agree almost everywhere, hence they agree at every point of \([0,1]\); the same is true of the differentiated identity \((V^\ast V f)^{\prime} = -Vf\). So we may evaluate these identities pointwise.

Evaluating \(V^\ast V f = \alpha f\) at \(x = 1\) gives \(\alpha f(1) = \int_1^1 (Vf) = 0\), so

\begin{equation*} f(1) = 0 . \end{equation*}

Evaluating \((V^\ast V f)^{\prime} = -Vf\) at \(x = 0\) gives \(\alpha f^{\prime}(0) = -(Vf)(0) = 0\), so

\begin{equation*} f^{\prime}(0) = 0 . \end{equation*}

Step 4: solving the eigenvalue problem.

Write \(\alpha = 1/\lambda^2\) where \(\lambda = \alpha^{-1/2} > 0\). Then Step 2 says \(f^{\prime\prime} + \lambda^2 f = 0\), whose solution space is the two-dimensional space spanned by \(x \mapsto \cos(\lambda x)\) and \(x \mapsto \sin(\lambda x)\). (Indeed, if \(f^{\prime\prime} + \lambda^2 f = 0\), then setting \(A = f(0)\) and \(B = f^{\prime}(0)/\lambda\), the function \(u = f - A\cos(\lambda \cdot) - B \sin(\lambda \cdot)\) satisfies \(u^{\prime\prime} = -\lambda^2 u\) and \(u(0) = u^{\prime}(0) = 0\). If \(u\) is real valued, then \((u^{\prime})^2 + \lambda^2 u^2\) has derivative \(2u^{\prime}(u^{\prime\prime} + \lambda^2 u) = 0\) and vanishes at \(0\), forcing \(u \equiv 0\); in the case \(\mathbb{F} = \mathbb{C}\), apply this to the real and imaginary parts of \(u\) separately, each of which satisfies the same real second-order equation with vanishing initial data.) Thus

\begin{equation*} f(x) = A \cos(\lambda x) + B \sin(\lambda x) \end{equation*}

for some scalars \(A, B\). The condition \(f^{\prime}(0) = 0\) forces \(B\lambda = 0\), hence \(B = 0\); since \(f \ne 0\) we get \(A \ne 0\). The condition \(f(1) = 0\) then forces

\begin{equation*} \cos \lambda = 0, \qquad \text{that is,} \qquad \lambda = \frac{(2k+1)\pi}{2} \text{ for some } k \in \{0, 1, 2, \dots\} \end{equation*}

(the positivity of \(\lambda\) rules out the negative odd multiples of \(\pi/2\)). Hence

\begin{equation*} \alpha = \frac{1}{\lambda^2} = \frac{4}{(2k+1)^2 \pi^2}, \end{equation*}

and the eigenspace corresponding to this \(\alpha\) is contained in the one-dimensional space spanned by \(x \mapsto \cos\!\big(\tfrac{(2k+1)\pi x}{2}\big)\).

Step 5: these numbers really are eigenvalues, with multiplicity \(1\).

Fix \(k \in \{0,1,2,\dots\}\), put \(\lambda_k = \frac{(2k+1)\pi}{2}\), and define \(f_k(x) = \cos(\lambda_k x)\). Then

\begin{equation*} (Vf_k)(x) = \int_0^x \cos(\lambda_k t) \, dt = \frac{\sin(\lambda_k x)}{\lambda_k}, \end{equation*}

and therefore

\begin{equation*} (V^\ast V f_k)(x) = \int_x^1 \frac{\sin(\lambda_k t)}{\lambda_k} \, dt = \frac{\cos(\lambda_k x) - \cos \lambda_k}{\lambda_k^{\,2}} = \frac{1}{\lambda_k^{\,2}} f_k(x), \end{equation*}

where we used \(\cos \lambda_k = \cos\!\big(\tfrac{(2k+1)\pi}{2}\big) = 0\). Thus \(f_k\) is an eigenvector of \(V^\ast V\) with eigenvalue \(1/\lambda_k^{\,2} = \frac{4}{(2k+1)^2\pi^2}\).

By Step 4 the corresponding eigenspace is exactly the one-dimensional span of \(f_k\), so each of these eigenvalues has geometric multiplicity \(1\). Moreover the numbers \(\lambda_0 < \lambda_1 < \lambda_2 < \cdots\) are distinct, so the eigenvalues \(1/\lambda_k^{\,2}\) are distinct and Step 4 shows there are no other positive eigenvalues of \(V^\ast V\).

Step 6: the singular values.

Taking positive square roots and listing each with its multiplicity \(1\) in decreasing order, Definition 10.116 gives

\begin{equation*} s_n(V) = \frac{1}{\lambda_{n-1}} = \frac{2}{(2n-1)\pi} \qquad \text{for } n \in \mathbb{Z}^+, \end{equation*}

so the list of singular values of the Volterra operator is

\begin{equation*} \frac{2}{\pi} \ge \frac{2}{3\pi} \ge \frac{2}{5\pi} \ge \cdots . \end{equation*}

In particular, \(V^\ast V\) has infinitely many positive eigenvalues, so the second bullet point of 10.116 never applies and no singular value of \(V\) equals \(0\); and \(s_n(V) \to 0\) as \(n \to \infty\), in accordance with 10.93 applied to the compact operator \(V^\ast V\).

Combining this with Exercise 12, which states that \(s_1(T) = \|T\|\) for every compact operator \(T\), we conclude

\begin{equation*} \|V\| = s_1(V) = \frac{2}{\pi} . \end{equation*}

Remark 1 (an explicit singular value decomposition).

Since \(\int_0^1 \cos^2(\lambda_k x) \, dx = \frac12 + \frac{\sin(2\lambda_k)}{4\lambda_k} = \frac12\) (because \(\sin(2\lambda_k) = \sin((2k+1)\pi) = 0\)), the functions

\begin{equation*} e_k(x) = \sqrt{2} \cos\!\Big( \frac{(2k+1)\pi x}{2} \Big), \qquad k = 0, 1, 2, \dots, \end{equation*}

form an orthonormal family of eigenvectors of \(V^\ast V\) (orthogonality holds because eigenvectors of a normal operator corresponding to distinct eigenvalues are orthogonal, by 10.57). This family is in fact an orthonormal basis of \(L^2([0,1])\): the operator \(V\) is injective (as noted in the paragraph after the proof of 10.70), so \(\langle V^\ast V f, f\rangle = \|Vf\|^2 = 0\) forces \(f = 0\) and thus \(V^\ast V\) is injective; hence, in the orthonormal basis of eigenvectors of \(V^\ast V\) provided by 10.106(a), every basis vector has a nonzero (necessarily positive) eigenvalue, and by Step 4 each such basis vector is a scalar multiple of some \(e_k\). Following the proof of 10.113, the second family is \(h_k = V e_k / s_{k+1}\), namely

\begin{equation*} h_k(x) = \sqrt{2} \sin\!\Big( \frac{(2k+1)\pi x}{2} \Big), \end{equation*}

and the singular value decomposition of the Volterra operator reads

\begin{equation*} Vf = \sum_{k=0}^{\infty} \frac{2}{(2k+1)\pi} \, \langle f, e_k \rangle \, h_k \qquad \text{for all } f \in L^2([0,1]). \end{equation*}

Remark 2 (a consistency check via 10.120).

With \(K\) as in Step 1 we have

\begin{equation*} \|K\|_{L^2(\mu \times \mu)}^2 = \int_0^1 \int_0^1 \chi_{\{y<x\}} \, dy \, dx = \int_0^1 x \, dx = \frac12 , \end{equation*}

while

\begin{equation*} \sum_{n=1}^{\infty} s_n(V)^2 = \sum_{k=0}^{\infty} \frac{4}{(2k+1)^2 \pi^2} = \frac{4}{\pi^2} \cdot \frac{\pi^2}{8} = \frac12 , \end{equation*}

using Example 10.124. The two agree, as 10.120 requires.

Remark 3 (comparison with Example 10.118).

The singular values of \(V\) are the same numbers \(\frac{2}{(2k+1)\pi}\) that appear as the singular values of \(V - V^\ast\) in Example 10.118, but there each value occurs twice (the corresponding eigenvalue of \((V-V^\ast)^\ast(V-V^\ast)\) has geometric multiplicity \(2\), coming from the two integers \(-k-1\) and \(k\)), whereas for \(V\) itself each occurs once. (Example 10.118 works in the case \(\mathbb{F} = \mathbb{C}\), as does Example 10.108 on which it depends.) Accordingly \(\|V - V^\ast\| = \frac{2}{\pi} = \|V\|\), even though \(\|V - V^\ast\| \le \|V\| + \|V^\ast\| = \frac{4}{\pi}\) (using \(\|V^\ast\| = \|V\|\), by 10.11) is all that the triangle inequality provides.

Fourier Analysis

Exercises 11A

Problem (11A.1)

Prove that \(\widehat{\overline{f}}(n) = \overline{\widehat{f}(-n)}\) for all \(f \in L^1(\partial \mathbf{D})\) and all \(n \in \mathbf{Z}\).

Solution

Suppose \(f \in L^1(\partial \mathbf{D})\) and \(n \in \mathbf{Z}\). Because \(|\overline{f}| = |f|\), the function \(\overline{f}\) also belongs to \(L^1(\partial \mathbf{D})\), so \(\widehat{\overline{f}}(n)\) is defined.

Recall from 11.7 that

\begin{equation*} \widehat{g}(n) = \int_{\partial \mathbf{D}} g(z)\, \overline{z^n} \, d\sigma(z) \end{equation*}

for \(g \in L^1(\partial \mathbf{D})\) and \(n \in \mathbf{Z}\).

The key pointwise identity is the following. If \(z \in \partial \mathbf{D}\), then \(|z| = 1\), hence \(z\overline{z} = 1\) and thus \(\overline{z} = 1/z\). Therefore

\begin{equation*} \overline{z^{\,n}} = \overline{z}^{\,n} = \frac{1}{z^{n}} = z^{-n} \end{equation*}

for every \(n \in \mathbf{Z}\). Consequently \(\overline{z^{-n}} = z^{n}\) as well, and so, using that complex conjugation is multiplicative,

\begin{equation*} \overline{f(z)}\;\overline{z^{\,n}} \;=\; \overline{f(z)}\cdot z^{-n} \;=\; \overline{\,f(z) \cdot \overline{z^{-n}}\,} \end{equation*}

for all \(z \in \partial \mathbf{D}\). The second equality holds because conjugation is an involution and is multiplicative, so that \(\overline{\,f(z)\,\overline{z^{-n}}\,} = \overline{f(z)} \cdot \overline{\overline{z^{-n}}} = \overline{f(z)}\, z^{-n}\).

Now integrate. Complex conjugation commutes with integration: if \(g \in L^1(\sigma)\), then writing \(g = (\operatorname{Re} g) + i(\operatorname{Im} g)\) and using the definition of the integral of a complex-valued function gives

\begin{equation*} \int \overline{g} \, d\sigma = \int (\operatorname{Re} g)\, d\sigma - i \int (\operatorname{Im} g) \, d\sigma = \overline{\int g \, d\sigma}. \end{equation*}

Applying this with \(g(z) = f(z)\,\overline{z^{-n}}\), which is in \(L^1(\sigma)\) because \(|g| = |f|\) on \(\partial \mathbf{D}\), we obtain

\begin{equation*} \begin{aligned} \widehat{\overline{f}}(n) &= \int_{\partial \mathbf{D}} \overline{f(z)}\; \overline{z^{\,n}} \, d\sigma(z) \\ &= \int_{\partial \mathbf{D}} \overline{\,f(z)\,\overline{z^{-n}}\,} \, d\sigma(z) \\ &= \overline{\int_{\partial \mathbf{D}} f(z)\, \overline{z^{-n}} \, d\sigma(z)} \\ &= \overline{\widehat{f}(-n)}, \end{aligned} \end{equation*}

as desired.

For an equivalent computation in the coordinates of \((-\pi, \pi]\): by 11.7,

\begin{equation*} \widehat{\overline{f}}(n) = \int_{-\pi}^{\pi} \overline{f(e^{it})}\, e^{-int} \, \frac{dt}{2\pi} = \overline{\int_{-\pi}^{\pi} f(e^{it})\, e^{int} \, \frac{dt}{2\pi}} = \overline{\widehat{f}(-n)}, \end{equation*}

since \(\overline{e^{-int}} = e^{int}\) and \(e^{int} = e^{-i(-n)t}\).

Problem (11A.2)

Suppose \(1 \le p \le \infty\) and \(n \in \mathbf{Z}\).

(a)
Show that the function \(f \mapsto \widehat{f}(n)\) is a bounded linear functional on \(L^p(\partial \mathbf{D})\) with norm \(1\).
(b)
Find all \(f \in L^p(\partial \mathbf{D})\) such that \(\|f\|_p = 1\) and \(|\widehat{f}(n)| = 1\).
Solution

Throughout, write \(\varphi(f) = \widehat{f}(n)\).

First note that \(\varphi\) is defined on all of \(L^p(\partial \mathbf{D})\), because \(\widehat{f}(n)\) is defined for \(f \in L^1(\partial \mathbf{D})\) and \(L^p(\partial \mathbf{D}) \subseteq L^1(\partial \mathbf{D})\) here. Indeed \(\sigma(\partial \mathbf{D}) = 1 < \infty\). The case \(p = 1\) is trivial. If \(1 < p < \infty\), then 7.10 (applied with the two exponents \(1 < p\), both finite) gives \(L^p(\partial \mathbf{D}) \subseteq L^1(\partial \mathbf{D})\) together with

\begin{equation*} \|f\|_1 \le \sigma(\partial \mathbf{D})^{(p-1)/p} \|f\|_p = \|f\|_p ; \end{equation*}

if \(p = \infty\) then \(\|f\|_1 = \int |f| \, d\sigma \le \|f\|_\infty \sigma(\partial \mathbf{D}) = \|f\|_\infty\), so again \(L^\infty(\partial \mathbf{D}) \subseteq L^1(\partial \mathbf{D})\) and \(\|f\|_1 \le \|f\|_p\). Thus in all cases

\begin{equation*} \|f\|_1 \le \|f\|_p \quad \text{for all } f \in L^p(\partial \mathbf{D}). \end{equation*}

Call this inequality \((*)\).

(a) Linearity of \(\varphi\) is 11.9(a) and 11.9(b). By 11.9(c) and \((*)\),

\begin{equation*} |\varphi(f)| = |\widehat{f}(n)| \le \|f\|_1 \le \|f\|_p \end{equation*}

for every \(f \in L^p(\partial \mathbf{D})\). Hence \(\varphi\) is a bounded linear functional on \(L^p(\partial \mathbf{D})\) with \(\|\varphi\| \le 1\).

For the reverse inequality, let \(g\) be the function \(z \mapsto z^n\) on \(\partial \mathbf{D}\). Then \(|g(z)| = 1\) for all \(z \in \partial \mathbf{D}\), so \(\|g\|_p = 1\) for every \(p \in [1, \infty]\) (using \(\sigma(\partial \mathbf{D}) = 1\)). Moreover, by 11.6,

\begin{equation*} \varphi(g) = \widehat{g}(n) = \int_{\partial \mathbf{D}} z^n \, \overline{z^n} \, d\sigma(z) = \int_{\partial \mathbf{D}} 1 \, d\sigma = 1 . \end{equation*}

Thus \(\|\varphi\| \ge |\varphi(g)| / \|g\|_p = 1\). Combining the two inequalities gives \(\|\varphi\| = 1\).

(b) We claim:

  • If \(1 < p \le \infty\), the functions in question are exactly \(f(z) = \lambda z^n\) with \(\lambda \in \mathbf{C}\), \(|\lambda| = 1\).
  • If \(p = 1\), they are exactly the functions of the form \(f(z) = \lambda h(z) z^n\), where \(\lambda \in \mathbf{C}\) with \(|\lambda| = 1\) and \(h \in L^1(\partial \mathbf{D})\) satisfies \(h \ge 0\) almost everywhere and \(\int_{\partial \mathbf{D}} h \, d\sigma = 1\). (Equivalently: \(\|f\|_1 = 1\) and the function \(z \mapsto f(z)\overline{z^n}\) has constant argument almost everywhere on the set where it is nonzero.)

Note that the \(p = 1\) family really is larger: for instance \(n = 0\) and \(h = 2 \cdot \chi_E\) with \(\sigma(E) = \tfrac12\) gives \(\|f\|_1 = 1\) and \(\widehat{f}(0) = 1\), while \(|f|\) is not constant.

Verification that these functions work: if \(|\lambda| = 1\), \(h \ge 0\), and \(\int h \, d\sigma = 1\), put \(f(z) = \lambda h(z) z^n\). Then \(\|f\|_1 = \int |h| \, d\sigma = 1\) and

\begin{equation*} \widehat{f}(n) = \int_{\partial \mathbf{D}} \lambda h(z) z^n \overline{z^n} \, d\sigma(z) = \lambda \int_{\partial \mathbf{D}} h \, d\sigma = \lambda, \end{equation*}

so \(|\widehat{f}(n)| = 1\). Taking \(h = 1\) recovers \(f(z) = \lambda z^n\), which has \(\|f\|_p = 1\) for every \(p \in [1, \infty]\) and \(\widehat{f}(n) = \lambda\).

Verification that there are no others: suppose \(f \in L^p(\partial \mathbf{D})\), \(\|f\|_p = 1\), and \(|\widehat{f}(n)| = 1\). By 11.9(c) and \((*)\),

\begin{equation*} 1 = |\widehat{f}(n)| \le \int_{\partial \mathbf{D}} |f| \, d\sigma = \|f\|_1 \le \|f\|_p = 1 , \end{equation*}

so both inequalities above are equalities.

Step 1: the equality \(|\widehat{f}(n)| = \int |f| \, d\sigma\). Set \(g(z) = f(z) \overline{z^n}\), so \(|g| = |f|\) and \(\widehat{f}(n) = \int g \, d\sigma\). Choose \(\lambda \in \mathbf{C}\) with \(|\lambda| = 1\) and \(\overline{\lambda} \int g \, d\sigma = \left| \int g \, d\sigma \right|\). Then

\begin{equation*} \int_{\partial \mathbf{D}} \bigl( |g| - \operatorname{Re}(\overline{\lambda} g) \bigr) d\sigma = \int |g| \, d\sigma - \operatorname{Re}\Bigl( \overline{\lambda} \int g \, d\sigma \Bigr) = \|f\|_1 - |\widehat{f}(n)| = 0 . \end{equation*}

The integrand is nonnegative (because \(\operatorname{Re} w \le |w|\) for all \(w \in \mathbf{C}\)) and integrates to \(0\); by 3.43 a nonnegative function in \(L^1(\sigma)\) with integral \(0\) vanishes almost everywhere, so \(\operatorname{Re}(\overline{\lambda} g) = |g|\) almost everywhere. Since \(|\overline{\lambda} g| = |g|\), this forces \(\overline{\lambda} g = |g|\) almost everywhere; that is,

\begin{equation*} f(z) \overline{z^n} = \lambda |f(z)| \qquad \text{for almost every } z \in \partial \mathbf{D}. \end{equation*}

Multiplying by \(z^n\) and using \(\overline{z^n} z^n = 1\) on \(\partial \mathbf{D}\) gives

\begin{equation*} f(z) = \lambda |f(z)| z^n \qquad \text{for almost every } z \in \partial \mathbf{D}. \end{equation*}

Call this equation \((**)\). Since \(\|f\|_1 = 1\), taking \(h = |f|\) shows \(f\) has exactly the form claimed for \(p = 1\). This completes the case \(p = 1\).

Step 2: the equality \(\|f\|_1 = \|f\|_p\) when \(p > 1\). Suppose first \(1 < p < \infty\) and put \(g = |f|\), so \(g \ge 0\), \(\int g \, d\sigma = 1\), and \(\int g^p \, d\sigma = 1\). For \(t \ge 0\) let \(\psi(t) = t^p - pt + p - 1\). Then \(\psi^{\prime}(t) = p(t^{p-1} - 1)\), which is negative on \((0,1)\) and positive on \((1, \infty)\); hence \(\psi\) has a strict global minimum at \(t = 1\), where \(\psi(1) = 0\). Thus \(\psi \ge 0\) on \([0, \infty)\), with \(\psi(t) = 0\) only for \(t = 1\). Now

\begin{equation*} \int_{\partial \mathbf{D}} \psi(g) \, d\sigma = \int g^p \, d\sigma - p \int g \, d\sigma + (p-1)\sigma(\partial \mathbf{D}) = 1 - p + p - 1 = 0 . \end{equation*}

Because \(\psi \circ g \ge 0\) and its integral is \(0\), we get \(\psi(g) = 0\) almost everywhere, hence \(g = 1\) almost everywhere; that is, \(|f| = 1\) almost everywhere.

If \(p = \infty\), then \(|f| \le \|f\|_\infty = 1\) almost everywhere and \(\int |f| \, d\sigma = 1 = \sigma(\partial \mathbf{D})\), so \(\int (1 - |f|) \, d\sigma = 0\) with nonnegative integrand, giving again \(|f| = 1\) almost everywhere.

In either case, substituting \(|f| = 1\) into \((**)\) yields \(f(z) = \lambda z^n\) for almost every \(z \in \partial \mathbf{D}\), with \(|\lambda| = 1\). This completes the proof.

Problem (11A.3)

Show that if \(0 \le r < 1\) and \(t \in \mathbf{R}\), then

\begin{equation*} P_r(e^{it}) = \frac{1 - r^2}{1 - 2r\cos t + r^2}. \end{equation*}

Solution

By 11.14, the Poisson kernel is defined by

\begin{equation*} P_r(\zeta) = \frac{1 - r^2}{|1 - r\zeta|^2} \end{equation*}

for \(\zeta \in \partial \mathbf{D}\) and \(r \in [0,1)\). So all that must be done is to evaluate \(|1 - re^{it}|^2\).

For any \(w \in \mathbf{C}\) we have \(|w|^2 = w \overline{w}\). Because \(r\) is real, \(\overline{1 - re^{it}} = 1 - re^{-it}\). Hence

\begin{equation*} \begin{aligned} |1 - re^{it}|^2 &= (1 - re^{it})(1 - re^{-it}) \\ &= 1 - re^{-it} - re^{it} + r^2 e^{it} e^{-it} \\ &= 1 - r(e^{it} + e^{-it}) + r^2 \\ &= 1 - 2r\cos t + r^2, \end{aligned} \end{equation*}

where the last equality uses the identity \(e^{it} + e^{-it} = 2\cos t\).

Therefore

\begin{equation*} P_r(e^{it}) = \frac{1 - r^2}{|1 - re^{it}|^2} = \frac{1 - r^2}{1 - 2r\cos t + r^2}, \end{equation*}

as desired.

Two remarks. First, the denominator never vanishes: \(1 - 2r\cos t + r^2 = (1-r)^2 + 2r(1 - \cos t) \ge (1-r)^2 > 0\) because \(r < 1\) and \(\cos t \le 1\). This rewriting also makes visible the two features of the Poisson kernel used repeatedly in this section: for fixed \(r\) close to \(1\) the denominator is small (of order \((1-r)^2\)) only when \(t\) is near \(0\), while for \(|t| \ge \delta\) we have \(1 - 2r\cos t + r^2 \ge 2r(1 - \cos\delta)\), so that \(P_r(e^{it}) \le \frac{1-r^2}{2r(1-\cos\delta)} \to 0\) as \(r \uparrow 1\), uniformly in such \(t\). Second, the formula shows that \(t \mapsto P_r(e^{it})\) is an even function of \(t\), a symmetry used in Exercise 5.

Problem (11A.4)

Suppose \(f \in L^1(\partial \mathbf{D})\), \(z \in \partial \mathbf{D}\), and \(f\) is continuous at \(z\). Prove that

\begin{equation*} \lim_{r \uparrow 1} (\mathcal{P}_r f)(z) = f(z). \end{equation*}

[The result in this exercise differs from 11.18 because here we are assuming continuity only at a single point and we are not even assuming that \(f\) is bounded, as compared to 11.18, which assumed continuity at all points of \(\partial \mathbf{D}\).]

Solution

Fix a representative of \(f\) that is defined at every point of \(\partial \mathbf{D}\) and is continuous at \(z\) (the hypothesis of continuity at \(z\) presupposes such a choice).

Two preliminary observations.

First, by 11.17 applied to \(h = P_r\) and by 11.16(b),

\begin{equation*} \int_{\partial \mathbf{D}} P_r(z\overline{w}) \, d\sigma(w) = \int_{\partial \mathbf{D}} P_r(\zeta) \, d\sigma(\zeta) = 1 \end{equation*}

for each \(r \in [0,1)\).

Second, if \(w \in \partial \mathbf{D}\), then \(w\overline{w} = 1\), so

\begin{equation*} 1 - z\overline{w} = \overline{w}\,(w - z) \quad \text{and hence} \quad |1 - z\overline{w}| = |w - z| . \end{equation*}

Now let \(\varepsilon > 0\). By continuity of \(f\) at \(z\), there exists \(\delta > 0\) such that

\begin{equation*} |f(w) - f(z)| < \varepsilon \quad \text{for all } w \in \partial \mathbf{D} \text{ with } |w - z| < \delta . \end{equation*}

Using the integral formula 11.15 and the first observation above,

\begin{equation*} (\mathcal{P}_r f)(z) - f(z) = \int_{\partial \mathbf{D}} \bigl( f(w) - f(z) \bigr) P_r(z\overline{w}) \, d\sigma(w) . \end{equation*}

All integrals here are defined: \(P_r\) is continuous on the compact set \(\partial \mathbf{D}\), hence bounded, and \(f \in L^1(\partial \mathbf{D})\) with \(\sigma(\partial \mathbf{D}) = 1\), so \(w \mapsto (f(w) - f(z))P_r(z\overline{w})\) is in \(L^1(\sigma)\).

Set \(A_\delta = \{w \in \partial \mathbf{D} : |w - z| < \delta\}\). This set is relatively open in \(\partial \mathbf{D}\), so \(\{t \in (-\pi, \pi] : e^{it} \in A_\delta\}\) is open in \((-\pi, \pi]\) and hence Borel; thus \(A_\delta\) is a measurable subset of \(\partial \mathbf{D}\) in the sense of 11.4. Splitting \(\partial \mathbf{D}\) into \(A_\delta\) and its complement and using \(P_r > 0\) (11.16(a)) gives

\begin{equation*} \bigl| (\mathcal{P}_r f)(z) - f(z) \bigr| \le \int_{A_\delta} |f(w) - f(z)| \, P_r(z\overline{w}) \, d\sigma(w) + \int_{\partial \mathbf{D} \setminus A_\delta} |f(w) - f(z)| \, P_r(z\overline{w}) \, d\sigma(w) . \end{equation*}

Bound the first integral: on \(A_\delta\) the integrand is at most \(\varepsilon P_r(z\overline{w})\), so this integral is at most

\begin{equation*} \varepsilon \int_{\partial \mathbf{D}} P_r(z\overline{w}) \, d\sigma(w) = \varepsilon . \end{equation*}

Bound the second integral. Here is exactly where the argument departs from the proof of 11.18: instead of pulling \(\|f\|_\infty\) out (which is unavailable, since \(f\) need not be bounded), we pull out a uniform bound on the Poisson kernel and leave \(f\) inside as an \(L^1\) function. Suppose \(w \notin A_\delta\), so \(|1 - z\overline{w}| = |w - z| \ge \delta\). If moreover \(1 - r < \frac{\delta}{2}\), then, since \(|z\overline{w}| = 1\),

\begin{equation*} |1 - r z\overline{w}| = \bigl| (1 - z\overline{w}) - (r-1)z\overline{w} \bigr| \ge |1 - z\overline{w}| - (1 - r) > \delta - \tfrac{\delta}{2} = \tfrac{\delta}{2}. \end{equation*}

Hence, by the definition 11.14 of the Poisson kernel,

\begin{equation*} P_r(z\overline{w}) = \frac{1 - r^2}{|1 - rz\overline{w}|^2} < \frac{4(1 - r^2)}{\delta^2} \end{equation*}

for all \(w \in \partial \mathbf{D} \setminus A_\delta\), whenever \(1 - r < \frac{\delta}{2}\). Therefore, for such \(r\),

\begin{equation*} \int_{\partial \mathbf{D} \setminus A_\delta} |f(w) - f(z)| \, P_r(z\overline{w}) \, d\sigma(w) \le \frac{4(1 - r^2)}{\delta^2} \int_{\partial \mathbf{D}} \bigl( |f(w)| + |f(z)| \bigr) d\sigma(w) = \frac{4(1 - r^2)}{\delta^2} \bigl( \|f\|_1 + |f(z)| \bigr), \end{equation*}

where we used \(\sigma(\partial \mathbf{D}) = 1\) and \(f \in L^1(\partial \mathbf{D})\).

Combining the two bounds: for all \(r \in [0,1)\) with \(1 - r < \frac{\delta}{2}\),

\begin{equation*} \bigl| (\mathcal{P}_r f)(z) - f(z) \bigr| \le \varepsilon + \frac{4(1 - r^2)}{\delta^2}\bigl( \|f\|_1 + |f(z)| \bigr). \end{equation*}

The quantity \(\|f\|_1 + |f(z)|\) is a finite constant not depending on \(r\), and \(1 - r^2 \to 0\) as \(r \uparrow 1\), while \(\delta\) depends only on \(\varepsilon\). Hence there exists \(r_0 \in [0,1)\) such that

\begin{equation*} \bigl| (\mathcal{P}_r f)(z) - f(z) \bigr| < 2\varepsilon \quad \text{for all } r \in (r_0, 1). \end{equation*}

Because \(\varepsilon > 0\) was arbitrary, \(\lim_{r \uparrow 1} (\mathcal{P}_r f)(z) = f(z)\).

Problem (11A.5)

Suppose \(a, b \in \mathbf{C}\), \(f \in L^1(\partial \mathbf{D})\), \(z \in \partial \mathbf{D}\), \(\displaystyle\lim_{t \downarrow 0} f(e^{it}z) = a\), and \(\displaystyle\lim_{t \uparrow 0} f(e^{it}z) = b\). Prove that

\begin{equation*} \lim_{r \uparrow 1} (\mathcal{P}_r f)(z) = \frac{a + b}{2}. \end{equation*}

[If \(a \ne b\), then \(f\) is said to have a jump discontinuity at \(z\).]

Solution

As in Exercise 4, the hypothesis about one-sided limits refers to a genuine function on \(\partial \mathbf{D}\), so fix a representative of \(f\) defined at every point of \(\partial \mathbf{D}\) for which the two stated limits hold. Changing \(f\) on a set of \(\sigma\)-measure \(0\) changes neither \(\widehat{f}(n)\) for any \(n\) nor, therefore, \(\mathcal{P}_r f\), so the conclusion does not depend on this choice.

Write

\begin{equation*} p_r(t) = P_r(e^{it}) = \frac{1 - r^2}{1 - 2r\cos t + r^2} \end{equation*}

for \(r \in [0,1)\) and \(t \in \mathbf{R}\), the last equality being Exercise 3. Since \(\cos\) is even, \(p_r\) is an even function of \(t\); in particular \(P_r(e^{-it}) = p_r(t)\).

Step 1: a rotated integral formula for \(\mathcal{P}_r f\) at \(z\).

By 11.15,

\begin{equation*} (\mathcal{P}_r f)(z) = \int_{\partial \mathbf{D}} f(w) P_r(z \overline{w}) \, d\sigma(w). \end{equation*}

Next we record the rotation invariance in the form we need: for every \(G \in L^1(\sigma)\),

\begin{equation*} \int_{\partial \mathbf{D}} G(zw) \, d\sigma(w) = \int_{\partial \mathbf{D}} G(w) \, d\sigma(w) . \end{equation*}

This follows from two applications of 11.17. First apply 11.17 to the function \(h\) defined by \(h(\zeta) = G(\overline{\zeta})\), with \(\overline{z}\) in place of \(z\) (note \(\overline{z} \in \partial \mathbf{D}\)): since \(\overline{\,\overline{z}\,\overline{w}\,} = zw\), this gives \(\int G(zw) \, d\sigma(w) = \int h(\overline{z}\,\overline{w}) \, d\sigma(w) = \int h(\zeta) \, d\sigma(\zeta) = \int G(\overline{\zeta}) \, d\sigma(\zeta)\). Then apply 11.17 to \(h = G\) with \(z\) replaced by \(1\) to get \(\int G(\overline{\zeta}) \, d\sigma(\zeta) = \int G(\zeta) \, d\sigma(\zeta)\). Applying this to \(G(w) = f(w) P_r(z\overline{w})\) and using \(z\overline{z} = 1\), which gives \(z\overline{zw} = \overline{w}\), we obtain

\begin{equation*} (\mathcal{P}_r f)(z) = \int_{\partial \mathbf{D}} f(zw) P_r(\overline{w}) \, d\sigma(w) = \int_{-\pi}^{\pi} f(e^{it}z) \, P_r(e^{-it}) \, \frac{dt}{2\pi} = \int_{-\pi}^{\pi} f(e^{it}z) \, p_r(t) \, \frac{dt}{2\pi}, \end{equation*}

where the middle equality is the transfer of integration on \(\partial \mathbf{D}\) to \((-\pi, \pi]\) recorded in 11.4. The same rotation invariance gives

\begin{equation*} \int_{-\pi}^{\pi} |f(e^{it}z)| \, \frac{dt}{2\pi} = \|f\|_1 < \infty, \end{equation*}

so \(t \mapsto f(e^{it}z)\) is in \(L^1\) of \((-\pi, \pi]\) with normalized Lebesgue measure, and all integrals below are defined.

Step 2: splitting the normalization.

By 11.16(b) and 11.4,

\begin{equation*} \int_{-\pi}^{\pi} p_r(t) \, \frac{dt}{2\pi} = \int_{\partial \mathbf{D}} P_r(\zeta) \, d\sigma(\zeta) = 1 , \end{equation*}

and since \(p_r\) is even, the two halves contribute equally:

\begin{equation*} \int_{0}^{\pi} p_r(t) \, \frac{dt}{2\pi} = \int_{-\pi}^{0} p_r(t) \, \frac{dt}{2\pi} = \frac{1}{2} . \end{equation*}

Consequently

\begin{equation*} (\mathcal{P}_r f)(z) - \frac{a+b}{2} = \int_{0}^{\pi} \bigl( f(e^{it}z) - a \bigr) p_r(t) \, \frac{dt}{2\pi} \; + \; \int_{-\pi}^{0} \bigl( f(e^{it}z) - b \bigr) p_r(t) \, \frac{dt}{2\pi}. \end{equation*}

This is the identity that makes the one-sided limits usable: the Poisson kernel splits its unit mass evenly between the two sides of \(z\), so the two one-sided limits are averaged with equal weights.

Step 3: the estimate.

Let \(\varepsilon > 0\). By the hypotheses on the one-sided limits, there exists \(\delta \in (0, \pi)\) such that

\begin{equation*} |f(e^{it}z) - a| < \varepsilon \text{ for } 0 < t < \delta, \qquad |f(e^{it}z) - b| < \varepsilon \text{ for } -\delta < t < 0 . \end{equation*}

Split each of the two integrals in Step 2 at \(\pm\delta\).

Near \(0\): since \(p_r > 0\),

\begin{equation*} \left| \int_{0}^{\delta} \bigl( f(e^{it}z) - a \bigr) p_r(t) \frac{dt}{2\pi} \right| \le \varepsilon \int_{0}^{\pi} p_r(t) \frac{dt}{2\pi} = \frac{\varepsilon}{2}, \end{equation*}

and likewise the integral over \((-\delta, 0)\) is at most \(\frac{\varepsilon}{2}\) in absolute value. Together these contribute at most \(\varepsilon\).

Away from \(0\): for \(\delta \le |t| \le \pi\) we have \(\cos t \le \cos \delta\), so

\begin{equation*} 1 - 2r\cos t + r^2 = (1-r)^2 + 2r(1 - \cos t) \ge 2r(1 - \cos\delta), \end{equation*}

and therefore, for \(\tfrac12 \le r < 1\),

\begin{equation*} p_r(t) \le \frac{1 - r^2}{2r(1 - \cos\delta)} \le \frac{1 - r^2}{1 - \cos\delta} \qquad (\delta \le |t| \le \pi). \end{equation*}

Note \(1 - \cos\delta > 0\) because \(0 < \delta < \pi\). Hence, for \(\tfrac12 \le r < 1\),

\begin{equation*} \left| \int_{\delta}^{\pi} \bigl( f(e^{it}z) - a \bigr) p_r(t) \frac{dt}{2\pi} \right| + \left| \int_{-\pi}^{-\delta} \bigl( f(e^{it}z) - b \bigr) p_r(t) \frac{dt}{2\pi} \right| \le \frac{1 - r^2}{1 - \cos\delta} \bigl( \|f\|_1 + |a| + |b| \bigr), \end{equation*}

where we used \(\int_{-\pi}^{\pi} |f(e^{it}z)| \frac{dt}{2\pi} = \|f\|_1\) and \(\int_{-\pi}^{\pi} \frac{dt}{2\pi} = 1\).

Combining, for all \(r \in [\tfrac12, 1)\),

\begin{equation*} \left| (\mathcal{P}_r f)(z) - \frac{a+b}{2} \right| \le \varepsilon + \frac{1 - r^2}{1 - \cos\delta} \bigl( \|f\|_1 + |a| + |b| \bigr). \end{equation*}

Here \(\delta\) depends only on \(\varepsilon\), and \(\|f\|_1 + |a| + |b|\) is a finite constant independent of \(r\), while \(1 - r^2 \to 0\) as \(r \uparrow 1\). Thus there exists \(r_0 \in [\tfrac12, 1)\) with

\begin{equation*} \left| (\mathcal{P}_r f)(z) - \frac{a+b}{2} \right| < 2\varepsilon \quad \text{for all } r \in (r_0, 1). \end{equation*}

Since \(\varepsilon > 0\) was arbitrary, \(\lim_{r \uparrow 1}(\mathcal{P}_r f)(z) = \frac{a+b}{2}\).

Taking \(a = b = f(z)\) recovers Exercise 4, since the two one-sided limits equalling \(f(z)\) is precisely continuity of \(f\) at \(z\) along \(\partial \mathbf{D}\).

Problem (11A.6)

Prove that for each \(p \in [1, \infty)\), there exists \(f \in L^1(\partial \mathbf{D})\) such that

\begin{equation*} \sum_{n = -\infty}^{\infty} |\widehat{f}(n)|^p = \infty . \end{equation*}

Solution

Fix \(p \in [1, \infty)\). Let \(\ell^p(\mathbf{Z}) = L^p(\nu)\), where \(\nu\) is counting measure on \(\mathbf{Z}\); thus \(\ell^p(\mathbf{Z})\) consists of the families \(c = (c_n)_{n \in \mathbf{Z}}\) of complex numbers with \(\|c\|_p = \bigl( \sum_{n} |c_n|^p \bigr)^{1/p} < \infty\), and \(\ell^p(\mathbf{Z})\) is a Banach space by 7.24. Also \(L^1(\partial \mathbf{D})\) is a Banach space by 7.24.

Suppose, for a contradiction, that

\begin{equation*} \sum_{n = -\infty}^{\infty} |\widehat{f}(n)|^p < \infty \quad \text{for every } f \in L^1(\partial \mathbf{D}). \end{equation*}

Step 1: the map \(f \mapsto \widehat{f}\) is then bounded from \(L^1(\partial \mathbf{D})\) to \(\ell^p(\mathbf{Z})\).

Under our supposition we may define \(T : L^1(\partial \mathbf{D}) \to \ell^p(\mathbf{Z})\) by \(Tf = (\widehat{f}(n))_{n \in \mathbf{Z}}\). This map is linear by 11.9(a) and 11.9(b). We show its graph is closed. Suppose \(f_1, f_2, \ldots\) is a sequence in \(L^1(\partial \mathbf{D})\) with \(\lim_k f_k = f\) in \(L^1(\partial \mathbf{D})\) and \(\lim_k Tf_k = c\) in \(\ell^p(\mathbf{Z})\). Fix \(n \in \mathbf{Z}\). By 11.9(a), 11.9(b), and 11.9(c),

\begin{equation*} |\widehat{f_k}(n) - \widehat{f}(n)| = |\widehat{(f_k - f)}(n)| \le \|f_k - f\|_1 \to 0 , \end{equation*}

so \(\widehat{f_k}(n) \to \widehat{f}(n)\). On the other hand, \(|\widehat{f_k}(n) - c_n| \le \|Tf_k - c\|_p \to 0\), so \(\widehat{f_k}(n) \to c_n\). Hence \(c_n = \widehat{f}(n)\) for every \(n\), that is, \(c = Tf\). Thus \(\operatorname{graph}(T)\) is closed, and the Closed Graph Theorem (6.85) implies that \(T\) is bounded: there exists \(C < \infty\) with

\begin{equation*} \Bigl( \sum_{n = -\infty}^{\infty} |\widehat{f}(n)|^p \Bigr)^{1/p} \le C \|f\|_1 \quad \text{for all } f \in L^1(\partial \mathbf{D}). \end{equation*}

Step 2: the Dirichlet kernel violates this inequality.

For \(N \in \mathbf{Z}^+\) define the continuous function \(D_N : \partial \mathbf{D} \to \mathbf{C}\) by

\begin{equation*} D_N(w) = \sum_{n = -N}^{N} w^n . \end{equation*}

By 11.6, \(\int_{\partial \mathbf{D}} w^n \overline{w^m} \, d\sigma(w)\) equals \(1\) if \(m = n\) and equals \(0\) if \(m \ne n\); hence

\begin{equation*} \widehat{D_N}(m) = \begin{cases} 1 & \text{if } |m| \le N, \\ 0 & \text{if } |m| > N, \end{cases} \end{equation*}

and therefore

\begin{equation*} \Bigl( \sum_{m = -\infty}^{\infty} |\widehat{D_N}(m)|^p \Bigr)^{1/p} = (2N+1)^{1/p} . \end{equation*}

Now we bound \(\|D_N\|_1\) from above. For \(t \in (-\pi, \pi]\) with \(t \ne 0\), summing the geometric series and then multiplying numerator and denominator by \(e^{-it/2}\) gives

\begin{equation*} D_N(e^{it}) = \sum_{n = -N}^{N} e^{int} = e^{-iNt} \, \frac{e^{i(2N+1)t} - 1}{e^{it} - 1} = \frac{e^{i(N + \frac12)t} - e^{-i(N + \frac12)t}}{e^{it/2} - e^{-it/2}} = \frac{\sin\bigl( (N + \frac12) t \bigr)}{\sin(t/2)} . \end{equation*}

Consequently \(|D_N(e^{it})| \le \frac{1}{|\sin(t/2)|}\); also \(|D_N(e^{it})| \le 2N+1\) directly from the definition. Because \(\sin\) is concave on \([0, \frac{\pi}{2}]\) we have \(\sin u \ge \frac{2u}{\pi}\) for \(u \in [0, \frac{\pi}{2}]\), so \(|\sin(t/2)| \ge \frac{|t|}{\pi}\) for \(0 < |t| \le \pi\), whence

\begin{equation*} |D_N(e^{it})| \le \frac{\pi}{|t|} \qquad (0 < |t| \le \pi). \end{equation*}

Using the bound \(2N+1\) on \(\{|t| \le \frac1N\}\) and the bound \(\frac{\pi}{|t|}\) on \(\{\frac1N \le |t| \le \pi\}\) (note \(\frac1N \le \pi\)),

\begin{equation*} \|D_N\|_1 = \int_{-\pi}^{\pi} |D_N(e^{it})| \, \frac{dt}{2\pi} \le \frac{1}{2\pi}\Bigl( \frac{2(2N+1)}{N} + \int_{\frac1N \le |t| \le \pi} \frac{\pi}{|t|} \, dt \Bigr) = \frac{1}{2\pi} \Bigl( 4 + \frac{2}{N} + 2\pi \ln(\pi N) \Bigr). \end{equation*}

Since \(\frac{1}{2\pi}(4 + \frac2N) \le \frac{3}{\pi} < 1\) and \(\ln \pi < 1.2\), we conclude

\begin{equation*} \|D_N\|_1 \le 3 + \ln N \qquad \text{for all } N \in \mathbf{Z}^+ . \end{equation*}

Step 3: the contradiction.

Applying the inequality of Step 1 to \(f = D_N\) and using Step 2,

\begin{equation*} (2N+1)^{1/p} \le C \|D_N\|_1 \le C (3 + \ln N) \end{equation*}

for every \(N \in \mathbf{Z}^+\). But \(p < \infty\), so \(\frac{(2N+1)^{1/p}}{3 + \ln N} \to \infty\) as \(N \to \infty\) (a positive power of \(N\) grows faster than \(\ln N\)). This contradicts the displayed inequality for large \(N\).

Hence our supposition was false: there exists \(f \in L^1(\partial \mathbf{D})\) with

\begin{equation*} \sum_{n = -\infty}^{\infty} |\widehat{f}(n)|^p = \infty , \end{equation*}

as desired.

Two remarks. First, the case \(p = 2\) shows the contrast with \(L^2(\partial \mathbf{D})\): there, Bessel’s inequality (8.57) together with 11.6 gives \(\sum_n |\widehat{f}(n)|^2 \le \|f\|_2^2 < \infty\). So the conclusion above says in particular that the Fourier coefficients of an \(L^1\) function need not be square summable, even though by the Riemann–Lebesgue Lemma (11.10) they always tend to \(0\). Second, since \(p = 1\) is allowed, there is an \(f \in L^1(\partial \mathbf{D})\) whose Fourier coefficients are not absolutely summable; contrast this with 11.29, where twice continuous differentiability forces absolute summability. The argument above is not constructive: it produces \(f\) from the Closed Graph Theorem (whose proof rests on Baire’s theorem) rather than exhibiting one explicitly.

Problem (11A.7)

Suppose \(\zeta \in \partial \mathbf{D}\). Show that the function

\begin{equation*} w \mapsto \frac{1 - |w|^2}{|1 - \overline{\zeta} w|^2} \end{equation*}

is harmonic on \(\mathbf{C} \setminus \{\zeta\}\) by finding an analytic function on \(\mathbf{C} \setminus \{\zeta\}\) whose real part is the function above.

Solution

Define \(g : \mathbf{C} \setminus \{\zeta\} \to \mathbf{C}\) by

\begin{equation*} g(w) = \frac{1 + \overline{\zeta} w}{1 - \overline{\zeta} w}. \end{equation*}

We claim that \(g\) is analytic on \(\mathbf{C} \setminus \{\zeta\}\) and that \(\operatorname{Re} g\) is the given function; harmonicity then follows from the first bullet point of 11.20.

Step 1: \(g\) is analytic on \(\mathbf{C} \setminus \{\zeta\}\).

Both \(w \mapsto 1 + \overline{\zeta}w\) and \(w \mapsto 1 - \overline{\zeta}w\) are polynomials in \(w\) (the number \(\overline{\zeta}\) is a constant, not a function of \(w\)), hence entire. Thus \(g\) is analytic wherever the denominator is nonzero. Now \(1 - \overline{\zeta}w = 0\) if and only if \(w = 1/\overline{\zeta}\). Because \(|\zeta| = 1\) we have \(\zeta \overline{\zeta} = 1\), so \(1/\overline{\zeta} = \zeta\). Hence the only singularity of \(g\) is at \(w = \zeta\), and \(g\) is analytic on \(\mathbf{C} \setminus \{\zeta\}\), exactly as required.

Step 2: computing \(\operatorname{Re} g\).

Let \(u = \overline{\zeta} w\), and suppose \(w \ne \zeta\), so \(u \ne 1\). Multiplying numerator and denominator by \(\overline{1 - u}\),

\begin{equation*} \frac{1 + u}{1 - u} = \frac{(1+u)\overline{(1-u)}}{|1-u|^2} = \frac{(1+u)(1 - \overline{u})}{|1-u|^2} = \frac{1 - \overline{u} + u - u\overline{u}}{|1-u|^2} = \frac{(1 - |u|^2) + (u - \overline{u})}{|1-u|^2}. \end{equation*}

The number \(u - \overline{u}\) is purely imaginary and \(1 - |u|^2\) and \(|1-u|^2\) are real, so taking real parts gives

\begin{equation*} \operatorname{Re} \frac{1+u}{1-u} = \frac{1 - |u|^2}{|1 - u|^2}. \end{equation*}

Finally, since \(|\overline{\zeta}| = |\zeta| = 1\), we have \(|u| = |\overline{\zeta} w| = |w|\). Substituting \(u = \overline{\zeta}w\) therefore yields

\begin{equation*} (\operatorname{Re} g)(w) = \frac{1 - |\overline{\zeta}w|^2}{|1 - \overline{\zeta}w|^2} = \frac{1 - |w|^2}{|1 - \overline{\zeta}w|^2} \end{equation*}

for all \(w \in \mathbf{C} \setminus \{\zeta\}\).

Step 3: conclusion.

By Steps 1 and 2, the function \(w \mapsto \frac{1 - |w|^2}{|1 - \overline{\zeta}w|^2}\) is the real part of the function \(g\), which is analytic on the open set \(\mathbf{C} \setminus \{\zeta\}\). By the first bullet point of 11.20, the real part of an analytic function on an open subset of \(\mathbf{C}\) is harmonic there. Hence the given function is harmonic on \(\mathbf{C} \setminus \{\zeta\}\).

For completeness, here is the reason the real part of an analytic function is harmonic. Write \(g = u + iv\) with \(u = \operatorname{Re} g\) and \(v = \operatorname{Im} g\) real valued. Analyticity gives the Cauchy–Riemann equations \(D_1 u = D_2 v\) and \(D_2 u = -D_1 v\), and \(u, v\) have continuous partial derivatives of all orders (an analytic function is infinitely differentiable). Hence

\begin{equation*} D_1(D_1 u) + D_2(D_2 u) = D_1(D_2 v) + D_2(-D_1 v) = D_1(D_2 v) - D_2(D_1 v) = 0, \end{equation*}

the last equality by equality of mixed partial derivatives for continuously twice differentiable functions. Thus \(\Delta u = 0\), so \(u\) is harmonic in the sense of 11.19.

Remark: taking \(w = rz\) with \(r \in [0,1)\) and \(z \in \partial \mathbf{D}\) gives \(\frac{1 - |w|^2}{|1 - \overline{\zeta}w|^2} = \frac{1 - r^2}{|1 - r \overline{\zeta} z|^2} = P_r(\overline{\zeta} z)\), so this exercise identifies the Poisson kernel, viewed as a function of the point \(w\) of the disk, as a harmonic function. This is the fact underlying 11.22, the harmonicity of the Poisson integral.

Problem (11A.8)

Suppose \(f : \partial\mathbb{D} \to \mathbb{R}\) is the function defined by

\begin{equation*} f(x, y) = x^4 y \end{equation*}

for \((x, y) \in \mathbb{R}^2\) with \(x^2 + y^2 = 1\). Find a polynomial \(u\) of two variables \(x, y\) such that \(u\) is harmonic on \(\mathbb{R}^2\) and \(u|_{\partial\mathbb{D}} = f\).

[Of course, \(u|_{\mathbb{D}}\) is the Poisson integral of \(f\). However, here you are asked to find an explicit formula for \(u\) in closed form, without involving or computing an integral. It may help to think of \(f\) as defined by \(f(z) = (\operatorname{Re} z)^4 (\operatorname{Im} z)\) for \(z \in \partial\mathbb{D}\).]

Solution

The answer is

\begin{equation*} u(x,y) = \tfrac{1}{16}\bigl(5x^4 y - 10 x^2 y^3 + y^5 + 9 x^2 y - 3 y^3 + 2 y\bigr), \end{equation*}

which is exactly \(u(w) = \tfrac{1}{16}\operatorname{Im}\bigl(w^5 + 3w^3 + 2w\bigr)\) under the usual identification \(w = x + yi\).

Derivation. Identify \(\mathbb{R}^2\) with \(\mathbb{C}\), so that \(f(z) = (\operatorname{Re} z)^4(\operatorname{Im} z)\) for \(z \in \partial\mathbb{D}\). For \(z \in \partial\mathbb{D}\) we have

\begin{equation*} \operatorname{Re} z = \frac{z + \overline{z}}{2}, \qquad \operatorname{Im} z = \frac{z - \overline{z}}{2i}, \qquad z\overline{z} = 1 . \end{equation*}

Hence

\begin{equation*} f(z) = \frac{(z+\overline{z})^4 (z - \overline{z})}{2^4 \cdot 2i} = \frac{(z+\overline{z})^4 (z-\overline{z})}{32 i}. \end{equation*}

Expanding by the binomial theorem and using \(z\overline{z} = 1\) to simplify the mixed terms,

\begin{equation*} (z + \overline{z})^4 = z^4 + 4 z^3\overline{z} + 6 z^2 \overline{z}^{\,2} + 4 z \overline{z}^{\,3} + \overline{z}^{\,4} = z^4 + 4z^2 + 6 + 4\overline{z}^{\,2} + \overline{z}^{\,4}. \end{equation*}

Multiplying by \(z - \overline{z}\) and again using \(z\overline{z} = 1\) (so that, for instance, \(z^4\overline{z} = z^3\) and \(z\overline{z}^{\,2} = \overline{z}\)) gives

\begin{equation*} \begin{aligned} (z+\overline{z})^4(z - \overline{z}) &= \bigl(z^5 + 4z^3 + 6z + 4\overline{z} + \overline{z}^{\,3}\bigr) - \bigl(z^3 + 4z + 6\overline{z} + 4\overline{z}^{\,3} + \overline{z}^{\,5}\bigr)\\ &= \bigl(z^5 - \overline{z}^{\,5}\bigr) + 3\bigl(z^3 - \overline{z}^{\,3}\bigr) + 2\bigl(z - \overline{z}\bigr)\\ &= 2i \operatorname{Im}\bigl(z^5 + 3z^3 + 2z\bigr). \end{aligned} \end{equation*}

Therefore

\begin{equation*} f(z) = \frac{2i \operatorname{Im}(z^5 + 3z^3 + 2z)}{32 i} = \tfrac{1}{16}\operatorname{Im}\bigl(z^5 + 3z^3 + 2z\bigr) \qquad \text{for } z \in \partial\mathbb{D}. \end{equation*}

Definition of \(u\) and verification. Let \(g : \mathbb{C} \to \mathbb{C}\) be the polynomial \(g(w) = \tfrac{1}{16}(w^5 + 3w^3 + 2w)\) and define

\begin{equation*} u = \operatorname{Im} g . \end{equation*}

Since \(g\) is a polynomial it is analytic on all of \(\mathbb{C}\), so \(u = \operatorname{Im} g\) is harmonic on \(\mathbb{C} = \mathbb{R}^2\) by the first bullet point of 11.20. The computation above shows \(u(z) = f(z)\) for every \(z \in \partial\mathbb{D}\), i.e. \(u|_{\partial\mathbb{D}} = f\).

Finally we write \(u\) out as a polynomial in \(x\) and \(y\). With \(w = x + yi\),

\begin{equation*} \operatorname{Im} w = y, \qquad \operatorname{Im}(w^3) = 3x^2 y - y^3, \qquad \operatorname{Im}(w^5) = 5x^4 y - 10 x^2 y^3 + y^5, \end{equation*}

as one sees by picking out the odd-index terms of the binomial expansions of \((x+yi)^3\) and \((x+yi)^5\). Hence

\begin{equation*} u(x,y) = \tfrac{1}{16}\Bigl(5x^4 y - 10x^2 y^3 + y^5\Bigr) + \tfrac{3}{16}\Bigl(3x^2 y - y^3\Bigr) + \tfrac{2}{16} y, \end{equation*}

which is the displayed formula. (One can also check \(u|_{\partial\mathbb{D}} = f\) directly: substituting \(y^2 = 1 - x^2\) into \(16u\) gives \(y\bigl[(2 - 3 + 1) + (9 + 3 - 10 - 2)x^2 + (5 + 10 + 1)x^4\bigr] = 16x^4 y\).)

Remark. The exponents \(\pm 1, \pm 3, \pm 5\) appearing here index exactly the nonzero Fourier coefficients of \(f\). Indeed the identity \(f(z) = \tfrac{1}{16}\operatorname{Im}(z^5 + 3z^3 + 2z)\) on \(\partial\mathbb{D}\) expands (using \(\operatorname{Im}w = \frac{w - \overline{w}}{2i}\) and \(\overline{z}^{\,n} = z^{-n}\)) into a finite linear combination of the functions \(z^n\), so by the orthonormality 11.6 its coefficients are the Fourier coefficients of \(f\). Equivalently, writing \(z = e^{it}\), the computation above says \(f(e^{it}) = \tfrac{1}{8}\sin t + \tfrac{3}{16}\sin 3t + \tfrac{1}{16}\sin 5t\), so \(\widehat{f}(\pm 1) = \mp \tfrac{i}{16}\), \(\widehat{f}(\pm 3) = \mp\tfrac{3i}{32}\), \(\widehat{f}(\pm 5) = \mp\tfrac{i}{32}\), and all other Fourier coefficients vanish. Thus by 11.11, for \(w = rz \in \mathbb{D}\) we get \((\mathcal{P}_r f)(z) = \sum_n r^{|n|}\widehat{f}(n)z^n = \tfrac{1}{16}\operatorname{Im}(w^5 + 3w^3 + 2w) = u(w)\), confirming that \(u|_{\mathbb{D}}\) is indeed the Poisson integral of \(f\) in the sense of 11.22. By 11.23 the restriction of \(u\) to \(\overline{\mathbb{D}}\) is therefore the solution of the Dirichlet problem 11.21 on \(\mathbb{D}\) for the boundary data \(f\). (Uniqueness of that solution follows from the maximum principle for harmonic functions, which is not proved in this book and is not needed here; the exercise only asks for one such \(u\).)

Problem (11A.9)

Find a formula (in closed form, not as an infinite sum) for \(\mathcal{P}_r f\), where \(f\) is the function in the second bullet point of Example 11.8.

Solution

The function in question is \(f : \partial\mathbb{D} \to \mathbb{R}\) defined by

\begin{equation*} f(z) = \frac{1}{|3 - z|^2}, \end{equation*}

and Example 11.8 computes its Fourier coefficients:

\begin{equation*} \widehat{f}(n) = \frac{1}{8}\cdot\frac{1}{3^{|n|}} \qquad \text{for all } n \in \mathbb{Z}. \end{equation*}

Claim. For every \(r \in [0,1)\) and every \(z \in \partial\mathbb{D}\),

\begin{equation*} (\mathcal{P}_r f)(z) = \frac{9 - r^2}{8\,|3 - rz|^2} = \frac{9 - r^2}{8\,(9 - 6r\operatorname{Re} z + r^2)} . \end{equation*}

Proof. Fix \(r \in [0,1)\) and \(z \in \partial\mathbb{D}\). By the definition 11.11 of \(\mathcal{P}_r f\) and the formula for \(\widehat{f}(n)\) above,

\begin{equation*} (\mathcal{P}_r f)(z) = \sum_{n=-\infty}^{\infty} r^{|n|}\widehat{f}(n) z^n = \frac{1}{8}\sum_{n=-\infty}^{\infty} \frac{r^{|n|}}{3^{|n|}}\, z^n = \frac{1}{8}\sum_{n=-\infty}^{\infty} \rho^{|n|} z^n, \end{equation*}

where \(\rho := r/3\). The series converges absolutely because \(\rho \in [0, \tfrac13)\) and \(|z^n| = 1\).

Now apply the summation formula 11.13, which states that \(\sum_{n=-\infty}^{\infty}\rho^{|n|}\zeta^{\,n} = \dfrac{1 - \rho^2}{|1 - \rho\zeta|^2}\) for \(\zeta \in \partial\mathbb{D}\) and \(\rho \in [0,1)\). (Its proof is just the sum of two geometric series: \(\sum_{n \geq 0}(\rho\zeta)^n + \sum_{n\geq 1}(\rho\overline{\zeta})^{\,n}\), using \(\zeta^{-n} = \overline{\zeta}^{\,n}\) for \(|\zeta| = 1\).) Taking \(\zeta = z\) gives

\begin{equation*} (\mathcal{P}_r f)(z) = \frac{1}{8}\cdot\frac{1 - \frac{r^2}{9}}{\bigl|1 - \frac{r}{3}z\bigr|^2} = \frac{1}{8}\cdot\frac{\frac{1}{9}(9 - r^2)}{\frac{1}{9}|3 - rz|^2} = \frac{9 - r^2}{8\,|3 - rz|^2}, \end{equation*}

where in the middle step we multiplied numerator and denominator by \(9\), using \(\bigl|1 - \tfrac{r}{3}z\bigr|^2 = \tfrac19|3 - rz|^2\). Expanding the denominator with \(|z| = 1\),

\begin{equation*} |3 - rz|^2 = (3 - rz)\overline{(3 - rz)} = 9 - 3r(z + \overline{z}) + r^2 |z|^2 = 9 - 6r\operatorname{Re}z + r^2, \end{equation*}

which gives the second displayed form and proves the claim. \(\blacksquare\)

Two checks. First, setting \(r = 1\) in the formula gives \(\dfrac{9-1}{8|3-z|^2} = \dfrac{1}{|3-z|^2} = f(z)\), so \((\mathcal{P}_r f)(z)\to f(z)\) as \(r\uparrow 1\); the convergence is in fact uniform on \(\partial\mathbb{D}\), as it must be by 11.18 since \(f\) is continuous (directly: \(|3 - rz| \geq 3 - r \geq 2\) for \(r\in[0,1]\) and \(z \in \partial\mathbb{D}\), so both denominators are bounded below by \(32\), and the numerators and denominators are jointly continuous in \((r,z)\) on the compact set \([0,1]\times\partial\mathbb{D}\)). Second, in the notation of 11.22, the Poisson integral \(u\) of \(f\) is given for \(w = rz \in \mathbb{D}\) by

\begin{equation*} u(w) = (\mathcal{P}_r f)(z) = \frac{9 - |w|^2}{8\,|3 - w|^2}, \end{equation*}

since \(r = |w|\) and \(rz = w\). Harmonicity of this \(u\) on \(\mathbb{D}\) (indeed on \(\mathbb{C}\setminus\{3\}\)) is guaranteed by 11.22, and can also be seen directly: for \(w \neq 3\),

\begin{equation*} \operatorname{Re}\frac{3 + w}{3 - w} = \frac{\operatorname{Re}\bigl[(3+w)(3 - \overline{w})\bigr]}{|3 - w|^2} = \frac{9 + 3\operatorname{Re}w - 3\operatorname{Re}w - |w|^2}{|3-w|^2} = \frac{9 - |w|^2}{|3 - w|^2}, \end{equation*}

so \(u = \operatorname{Re} g\) where \(g(w) = \tfrac18\cdot\frac{3+w}{3-w}\) is analytic on \(\mathbb{C}\setminus\{3\}\); now apply the first bullet point of 11.20.

Problem (11A.10)

Suppose \(f : \partial\mathbb{D} \to \mathbb{C}\) is three times continuously differentiable. Prove that

\begin{equation*} f^{[1]}(z) = i \sum_{n=-\infty}^{\infty} n\,\widehat{f}(n)\, z^n \end{equation*}

for all \(z \in \partial\mathbb{D}\).

Solution

Throughout, \(\widetilde{g}\) denotes the transfer of \(g : \partial\mathbb{D}\to\mathbb{C}\) to \(\mathbb{R}\), so \(\widetilde{g}(t) = g(e^{it})\), and \(g^{[k]}\) is defined by \(g^{[k]}(e^{it}) = \widetilde{g}^{\,(k)}(t)\), as in 11.24.

Step 1: \(f^{[1]}\) is twice continuously differentiable.

By hypothesis \(\widetilde{f}\) is three times differentiable everywhere on \(\mathbb{R}\) and \(\widetilde{f}^{\,(3)}\) is continuous. Because a function that is differentiable at every point of \(\mathbb{R}\) is continuous there, the existence of \(\widetilde{f}^{\,(2)}\) and \(\widetilde{f}^{\,(3)}\) everywhere makes \(\widetilde{f}^{\,(1)}\) and \(\widetilde{f}^{\,(2)}\) continuous as well. Thus \(f\) is \(k\) times continuously differentiable in the sense of 11.24 for \(k = 1, 2, 3\), and in particular \(f^{[1]}\), \(f^{[2]}\), \(f^{[3]}\) are all defined.

Directly from the two bullet points of 11.24 we have, for every \(t \in \mathbb{R}\),

\begin{equation*} \widetilde{f^{[1]}}(t) = f^{[1]}(e^{it}) = \widetilde{f}^{\,(1)}(t), \qquad\text{so}\qquad \widetilde{f^{[1]}} = \widetilde{f}^{\,(1)} . \end{equation*}

Consequently \(\bigl(\widetilde{f^{[1]}}\bigr)^{(2)} = \widetilde{f}^{\,(3)}\), which exists everywhere on \(\mathbb{R}\) and is continuous. Hence \(\widetilde{f^{[1]}}\) is twice differentiable everywhere with continuous second derivative, which by 11.24 means exactly that the function \(f^{[1]} : \partial\mathbb{D}\to\mathbb{C}\) is twice continuously differentiable.

Step 2: the Fourier coefficients of \(f^{[1]}\).

Because \(f\) is (at least) once continuously differentiable, 11.26 with \(k = 1\) gives

\begin{equation*} \widehat{f^{[1]}}(n) = i^1 n^1 \widehat{f}(n) = i\, n\, \widehat{f}(n) \end{equation*}

for every \(n \in \mathbb{Z}\).

Step 3: apply 11.27 to \(f^{[1]}\).

By Step 1 the function \(f^{[1]}\) satisfies the hypothesis of 11.27. Hence

\begin{equation*} f^{[1]}(z) = \sum_{n=-\infty}^{\infty} \widehat{f^{[1]}}(n)\, z^n \end{equation*}

for all \(z \in \partial\mathbb{D}\), with the partial sums \(\sum_{n=-K}^{M}\) converging uniformly on \(\partial\mathbb{D}\) as \(K, M \to \infty\). Moreover the proof of 11.27 (specifically 11.29 applied to \(f^{[1]}\)) shows that

\begin{equation*} \sum_{n=-\infty}^{\infty}\bigl|\widehat{f^{[1]}}(n)\bigr| = \sum_{n=-\infty}^{\infty} |n|\,|\widehat{f}(n)| < \infty, \end{equation*}

so the doubly infinite series converges absolutely; in particular it is an unambiguous unordered sum and the constant \(i\) may be factored out of it.

Substituting the formula from Step 2 into the last display of Step 3 now gives, for every \(z \in \partial\mathbb{D}\),

\begin{equation*} f^{[1]}(z) = \sum_{n=-\infty}^{\infty} i\, n\, \widehat{f}(n)\, z^n = i \sum_{n=-\infty}^{\infty} n\, \widehat{f}(n)\, z^n, \end{equation*}

as desired. \(\blacksquare\)

Remark. The role of the third derivative is only to put \(f^{[1]}\) inside the reach of 11.27, which was stated for twice continuously differentiable functions. Since (as noted in the text just before 11.27) that theorem in fact holds for continuously differentiable functions, the conclusion above is already true when \(f\) is merely twice continuously differentiable. Note also that this identity is exactly what one gets by formally differentiating the Fourier series \(f(z) = \sum_n \widehat{f}(n) z^n\) with respect to \(t\) along \(z = e^{it}\): since \(\frac{d}{dt}e^{int} = i n e^{int}\), term-by-term differentiation produces \(i\sum_n n\widehat{f}(n)z^n\), and the point of the exercise is that this formal operation is legitimate under the stated smoothness.

Problem (11A.11)

Let \(C(\partial\mathbb{D})\) denote the Banach space of continuous functions from \(\partial\mathbb{D}\) to \(\mathbb{C}\), with the supremum norm. For \(M \in \mathbb{Z}^+\), define a linear functional \(\varphi_M : C(\partial\mathbb{D}) \to \mathbb{C}\) by

\begin{equation*} \varphi_M(f) = \sum_{n=-M}^{M} \widehat{f}(n). \end{equation*}

Thus \(\varphi_M(f)\) is a partial sum of the Fourier series \(\sum_{n=-\infty}^{\infty}\widehat{f}(n) z^n\), evaluated at \(z = 1\).

(a) Show that

\begin{equation*} \varphi_M(f) = \int_{-\pi}^{\pi} f(e^{it}) \, \frac{\sin\bigl(M + \tfrac12\bigr)t}{\sin\frac t2}\,\frac{dt}{2\pi} \end{equation*}

for every \(f \in C(\partial\mathbb{D})\) and every \(M \in \mathbb{Z}^+\).

(b) Show that

\begin{equation*} \lim_{M\to\infty} \int_{-\pi}^{\pi} \left| \frac{\sin\bigl(M+\tfrac12\bigr)t}{\sin\frac t2} \right| \frac{dt}{2\pi} = \infty . \end{equation*}

(c) Show that \(\lim_{M\to\infty}\|\varphi_M\| = \infty\).

(d) Show that there exists \(f \in C(\partial\mathbb{D})\) such that \(\displaystyle\lim_{M\to\infty}\sum_{n=-M}^{M}\widehat{f}(n)\) does not exist (as an element of \(\mathbb{C}\)).

[Because the sum in (d) is a partial sum of the Fourier series evaluated at \(z = 1\), part (d) shows that the Fourier series of a continuous function on \(\partial\mathbb{D}\) need not converge pointwise on \(\partial\mathbb{D}\).

The family of functions (one for each \(M \in \mathbb{Z}^+\)) on \(\partial\mathbb{D}\) defined by

\begin{equation*} e^{it} \mapsto \frac{\sin\bigl(M+\tfrac12\bigr)t}{\sin\frac t2} \end{equation*}

is called the Dirichlet kernel.]

Solution

Throughout, for \(M \in \mathbb{Z}^+\) define \(D_M : \mathbb{R}\to\mathbb{R}\) by

\begin{equation*} D_M(t) = \sum_{n=-M}^{M} e^{int} . \end{equation*}

This is a finite sum of continuous \(2\pi\)-periodic functions, hence \(D_M\) is continuous and \(2\pi\)-periodic on \(\mathbb{R}\); it is real-valued because pairing the terms \(n\) and \(-n\) gives \(D_M(t) = 1 + 2\sum_{n=1}^{M}\cos nt\). Also set

\begin{equation*} L_M = \int_{-\pi}^{\pi} |D_M(t)| \,\frac{dt}{2\pi}. \end{equation*}

Step 0: the closed form for \(D_M\).

Suppose \(t \in [-\pi,\pi]\) with \(t \neq 0\); then \(e^{it}\neq 1\), and summing the geometric series with ratio \(e^{it}\) gives

\begin{equation*} D_M(t) = e^{-iMt}\sum_{k=0}^{2M} e^{ikt} = e^{-iMt}\,\frac{e^{i(2M+1)t} - 1}{e^{it}-1} = \frac{e^{i(M+1)t} - e^{-iMt}}{e^{it}-1}. \end{equation*}

Multiplying numerator and denominator by \(e^{-it/2}\) yields

\begin{equation*} D_M(t) = \frac{e^{i(M+\frac12)t} - e^{-i(M+\frac12)t}}{e^{it/2} - e^{-it/2}} = \frac{2i\sin\bigl(M+\tfrac12\bigr)t}{2i\sin\frac t2} = \frac{\sin\bigl(M+\tfrac12\bigr)t}{\sin\frac t2}. \end{equation*}

Thus \(D_M(t)\) equals the Dirichlet kernel for all \(t \in [-\pi,\pi]\setminus\{0\}\); at \(t = 0\) the right side is undefined while \(D_M(0) = 2M+1\), and indeed \(\lim_{t\to 0}\frac{\sin(M+\frac12)t}{\sin\frac t2} = \frac{M+\frac12}{\frac12} = 2M+1\), so \(D_M\) is precisely the continuous extension of the Dirichlet kernel to \([-\pi,\pi]\). Since \(\{0\}\) has measure \(0\), we may freely replace the Dirichlet kernel by \(D_M\) inside any integral over \([-\pi,\pi]\).

(a) Let \(f \in C(\partial\mathbb{D})\) and \(M \in \mathbb{Z}^+\). By the definition 11.7 of Fourier coefficients,

\begin{equation*} \varphi_M(f) = \sum_{n=-M}^{M}\widehat{f}(n) = \sum_{n=-M}^{M}\int_{-\pi}^{\pi} f(e^{it})e^{-int}\,\frac{dt}{2\pi} = \int_{-\pi}^{\pi} f(e^{it})\Bigl(\sum_{n=-M}^{M} e^{-int}\Bigr)\frac{dt}{2\pi}, \end{equation*}

where interchanging the finite sum with the integral is just linearity of the integral (each integrand is continuous on the compact interval \([-\pi,\pi]\), hence integrable). Replacing \(n\) by \(-n\) in the inner sum (a bijection of \(\{-M,\dots,M\}\) onto itself) gives \(\sum_{n=-M}^{M}e^{-int} = D_M(t)\). Hence, using Step 0,

\begin{equation*} \varphi_M(f) = \int_{-\pi}^{\pi} f(e^{it})\, D_M(t)\,\frac{dt}{2\pi} = \int_{-\pi}^{\pi} f(e^{it})\,\frac{\sin\bigl(M+\tfrac12\bigr)t}{\sin\frac t2}\,\frac{dt}{2\pi}, \end{equation*}

as required.

(b) By the last sentence of Step 0, the integrand \(\bigl|\sin(M+\tfrac12)t / \sin\tfrac t2\bigr|\) agrees with \(|D_M|\) on \([-\pi,\pi]\setminus\{0\}\), so the integral in (b) equals \(L_M\); thus we must show \(L_M \to \infty\). Because \(|\sin s| \leq |s|\) for all \(s \in \mathbb{R}\), we have \(\bigl|\sin\frac t2\bigr| \leq \frac{|t|}{2}\), and therefore, for \(t \in [-\pi,\pi]\setminus\{0\}\),

\begin{equation*} |D_M(t)| = \frac{\bigl|\sin\bigl(M+\tfrac12\bigr)t\bigr|}{\bigl|\sin\frac t2\bigr|} \geq \frac{2\bigl|\sin\bigl(M+\tfrac12\bigr)t\bigr|}{|t|}. \end{equation*}

Since the integrand is even, this gives

\begin{equation*} L_M \geq \frac{1}{2\pi}\int_{-\pi}^{\pi}\frac{2\bigl|\sin\bigl(M+\tfrac12\bigr)t\bigr|}{|t|}\,dt = \frac{2}{\pi}\int_0^{\pi}\frac{\bigl|\sin\bigl(M+\tfrac12\bigr)t\bigr|}{t}\,dt = \frac{2}{\pi}\int_0^{(M+\frac12)\pi}\frac{|\sin s|}{s}\,ds, \end{equation*}

the last equality by the change of variable \(s = \bigl(M+\tfrac12\bigr)t\) (which leaves \(\frac{dt}{t} = \frac{ds}{s}\) invariant). Since the integrand is nonnegative and \(\bigl(M+\tfrac12\bigr)\pi > M\pi\),

\begin{equation*} \int_0^{(M+\frac12)\pi}\frac{|\sin s|}{s}\,ds \geq \int_0^{M\pi}\frac{|\sin s|}{s}\,ds = \sum_{k=0}^{M-1}\int_{k\pi}^{(k+1)\pi}\frac{|\sin s|}{s}\,ds \geq \sum_{k=0}^{M-1}\frac{1}{(k+1)\pi}\int_{k\pi}^{(k+1)\pi}|\sin s|\,ds, \end{equation*}

where we bounded \(\frac1s \geq \frac{1}{(k+1)\pi}\) on \([k\pi,(k+1)\pi]\). Each of those integrals equals \(2\), so

\begin{equation*} L_M \geq \frac{2}{\pi}\cdot\sum_{k=0}^{M-1}\frac{2}{(k+1)\pi} = \frac{4}{\pi^2}\sum_{k=1}^{M}\frac{1}{k}. \end{equation*}

The harmonic series diverges, so \(L_M \to \infty\) as \(M \to \infty\), proving (b).

(c) We show \(\|\varphi_M\| = L_M\); combined with (b) this gives \(\lim_{M\to\infty}\|\varphi_M\| = \infty\).

First, by (a), for every \(f \in C(\partial\mathbb{D})\),

\begin{equation*} |\varphi_M(f)| \leq \int_{-\pi}^{\pi}\bigl|f(e^{it})\bigr|\,|D_M(t)|\,\frac{dt}{2\pi} \leq \|f\|_\infty L_M , \end{equation*}

so \(\varphi_M\) is a bounded linear functional with \(\|\varphi_M\| \leq L_M\).

For the reverse inequality, let \(\varepsilon > 0\) and define \(f_\varepsilon : \partial\mathbb{D}\to\mathbb{C}\) by

\begin{equation*} f_\varepsilon(e^{it}) = \frac{D_M(t)}{|D_M(t)| + \varepsilon}\qquad (t \in \mathbb{R}). \end{equation*}

This is well defined (the denominator is at least \(\varepsilon > 0\)) and unambiguous (\(D_M\) is \(2\pi\)-periodic), and \(f_\varepsilon\) is continuous on \(\partial\mathbb{D}\) because \(D_M\) is continuous on \(\mathbb{R}\) and \(t \mapsto e^{it}\) is a continuous open surjection of \(\mathbb{R}\) onto \(\partial\mathbb{D}\) that identifies exactly the points differing by multiples of \(2\pi\). Clearly \(\|f_\varepsilon\|_\infty \leq 1\). By (a),

\begin{equation*} \varphi_M(f_\varepsilon) = \int_{-\pi}^{\pi}\frac{D_M(t)^2}{|D_M(t)|+\varepsilon}\,\frac{dt}{2\pi}. \end{equation*}

For every real number \(a\) we have

\begin{equation*} \frac{a^2}{|a| + \varepsilon} = |a| - \frac{\varepsilon |a|}{|a| + \varepsilon} \geq |a| - \varepsilon , \end{equation*}

since \(\frac{|a|}{|a|+\varepsilon} \leq 1\). Applying this with \(a = D_M(t)\) and integrating,

\begin{equation*} \|\varphi_M\| \geq |\varphi_M(f_\varepsilon)| = \varphi_M(f_\varepsilon) \geq \int_{-\pi}^{\pi}\bigl(|D_M(t)| - \varepsilon\bigr)\frac{dt}{2\pi} = L_M - \varepsilon . \end{equation*}

As \(\varepsilon > 0\) was arbitrary, \(\|\varphi_M\| \geq L_M\). Hence \(\|\varphi_M\| = L_M\), and (b) gives \(\lim_{M\to\infty}\|\varphi_M\| = \infty\).

(d) Suppose, for a contradiction, that for every \(f \in C(\partial\mathbb{D})\) the limit \(\lim_{M\to\infty}\sum_{n=-M}^{M}\widehat{f}(n) = \lim_{M\to\infty}\varphi_M(f)\) exists in \(\mathbb{C}\). A convergent sequence of complex numbers is bounded, so this would give

\begin{equation*} \sup_{M \in \mathbb{Z}^+} |\varphi_M(f)| < \infty \qquad\text{for every } f \in C(\partial\mathbb{D}). \end{equation*}

Now \(C(\partial\mathbb{D})\) with the supremum norm is a Banach space (as stated in the exercise), \(\mathbb{C}\) is a normed vector space, and by (c) each \(\varphi_M\) is a bounded linear functional on \(C(\partial\mathbb{D})\). Hence the Principle of Uniform Boundedness (6.86), applied to the family \(\{\varphi_M : M \in \mathbb{Z}^+\}\), would give

\begin{equation*} \sup_{M\in\mathbb{Z}^+}\|\varphi_M\| < \infty, \end{equation*}

contradicting (c), which says \(\|\varphi_M\| \to \infty\).

Therefore there exists \(f \in C(\partial\mathbb{D})\) such that \(\lim_{M\to\infty}\sum_{n=-M}^{M}\widehat{f}(n)\) does not exist in \(\mathbb{C}\). \(\blacksquare\)

Remark. Since \(\sum_{n=-M}^{M}\widehat{f}(n)\) is the \(M\)-th symmetric partial sum of the Fourier series of \(f\) evaluated at \(z = 1\), part (d) exhibits a continuous function on \(\partial\mathbb{D}\) whose Fourier series fails to converge at the point \(1\). This is the example referred to in the text before 11.11 and again before 11.27, and it is why the Abel-summation device \(\mathcal{P}_r f\) (11.11) is needed: by 11.18, \(\mathcal{P}_r f \to f\) uniformly for every continuous \(f\). The proof above is nonconstructive; the quantity \(L_M = \|\varphi_M\|\) is called the \(M\)-th Lebesgue constant, and the estimate in (b) is sharp up to a constant, with \(L_M = \frac{4}{\pi^2}\log M + O(1)\).

Problem (11A.12)

Define \(f : \partial\mathbb{D}\to\mathbb{R}\) by

\begin{equation*} f(z) = \begin{cases} 1 & \text{if } \operatorname{Im} z > 0,\\ -1 & \text{if } \operatorname{Im} z < 0,\\ 0 & \text{if } \operatorname{Im} z = 0.\end{cases} \end{equation*}

(a) Show that if \(n \in \mathbb{Z}\), then

\begin{equation*} \widehat{f}(n) = \begin{cases} -\dfrac{2i}{n\pi} & \text{if } n \text{ is odd},\\[2mm] 0 & \text{if } n \text{ is even}.\end{cases} \end{equation*}

(b) Show that

\begin{equation*} (\mathcal{P}_r f)(z) = \frac{2}{\pi}\arctan\frac{2r\operatorname{Im} z}{1 - r^2} \end{equation*}

for every \(r \in [0,1)\) and every \(z \in \partial\mathbb{D}\).

(c) Verify that \(\lim_{r\uparrow 1}(\mathcal{P}_r f)(z) = f(z)\) for every \(z \in \partial\mathbb{D}\).

(d) Prove that \(\mathcal{P}_r f\) does not converge uniformly to \(f\) on \(\partial\mathbb{D}\) as \(r \uparrow 1\).

Solution

For \(t \in (-\pi,\pi]\) we have \(\operatorname{Im}(e^{it}) = \sin t\), which is positive exactly for \(t \in (0,\pi)\), negative exactly for \(t \in (-\pi,0)\), and zero exactly for \(t \in \{0,\pi\}\). Hence

\begin{equation*} f(e^{it}) = \begin{cases} 1 & \text{if } t \in (0,\pi),\\ -1 & \text{if } t \in (-\pi,0),\\ 0 & \text{if } t \in \{0,\pi\}.\end{cases} \end{equation*}

In particular \(f\) is measurable (the two sets above are Borel) and \(|f| \leq 1\), so \(f \in L^1(\partial\mathbb{D})\) and \(\mathcal{P}_r f\) is defined by 11.11.

(a) By 11.7 and the description of \(f(e^{it})\) above (the two-point set \(\{0,\pi\}\) has measure zero and may be ignored),

\begin{equation*} \widehat{f}(n) = \int_{-\pi}^{\pi} f(e^{it}) e^{-int}\,\frac{dt}{2\pi} = \frac{1}{2\pi}\left(\int_0^{\pi} e^{-int}\,dt - \int_{-\pi}^{0} e^{-int}\,dt\right). \end{equation*}

If \(n = 0\), both integrals equal \(\pi\), so \(\widehat{f}(0) = 0\); this agrees with the claim since \(0\) is even.

Now suppose \(n \in \mathbb{Z}\setminus\{0\}\). Then

\begin{equation*} \int_0^{\pi} e^{-int}\,dt = \frac{e^{-int}}{-in}\bigg|_{t=0}^{t=\pi} = \frac{e^{-in\pi} - 1}{-in} = \frac{1 - (-1)^n}{in}, \end{equation*}

using \(e^{-in\pi} = (-1)^n\), and similarly

\begin{equation*} \int_{-\pi}^{0} e^{-int}\,dt = \frac{e^{-int}}{-in}\bigg|_{t=-\pi}^{t=0} = \frac{1 - e^{in\pi}}{-in} = -\,\frac{1 - (-1)^n}{in}. \end{equation*}

Subtracting and dividing by \(2\pi\),

\begin{equation*} \widehat{f}(n) = \frac{1}{2\pi}\cdot\frac{2\bigl(1 - (-1)^n\bigr)}{in} = \frac{1 - (-1)^n}{i n \pi}. \end{equation*}

If \(n\) is even then \(1 - (-1)^n = 0\), so \(\widehat{f}(n) = 0\). If \(n\) is odd then \(1 - (-1)^n = 2\), so, using \(\frac1i = -i\),

\begin{equation*} \widehat{f}(n) = \frac{2}{i n \pi} = -\frac{2i}{n\pi}, \end{equation*}

which is the asserted formula.

(b) Fix \(r \in [0,1)\) and \(z \in \partial\mathbb{D}\), and put \(w = rz\), so \(w \in \mathbb{D}\), \(|w| = r\) and \(\operatorname{Im} w = r \operatorname{Im} z\). By 11.11 and part (a),

\begin{equation*} (\mathcal{P}_r f)(z) = \sum_{n=-\infty}^{\infty} r^{|n|}\widehat{f}(n) z^n = \sum_{n \text{ odd}} r^{|n|}\Bigl(-\frac{2i}{n\pi}\Bigr) z^n . \end{equation*}

This series converges absolutely (its terms are dominated by \(\frac{2}{\pi}r^{|n|}\), summable since \(r < 1\)), so we may group the terms \(n\) and \(-n\) for each odd \(n \geq 1\). Using \(z^{-n} = \overline{z}^{\,n}\) for \(|z| = 1\) and \(z^n - \overline{z}^{\,n} = 2i \operatorname{Im}(z^n)\), the pair contributes

\begin{equation*} -\frac{2i}{n\pi}r^n z^n + \frac{2i}{n\pi}r^n \overline{z}^{\,n} = -\frac{2i}{n\pi}r^n\bigl(z^n - \overline{z}^{\,n}\bigr) = -\frac{2i}{n\pi}r^n\cdot 2i \operatorname{Im}(z^n) = \frac{4}{n\pi}\,\operatorname{Im}\bigl((rz)^n\bigr). \end{equation*}

Writing the odd positive integers as \(n = 2k+1\), we obtain

\begin{equation*} (\mathcal{P}_r f)(z) = \frac{4}{\pi}\sum_{k=0}^{\infty}\frac{\operatorname{Im}\bigl(w^{2k+1}\bigr)}{2k+1} = \frac{4}{\pi}\operatorname{Im}\, g(w), \qquad\text{where } g(w) := \sum_{k=0}^{\infty}\frac{w^{2k+1}}{2k+1}. \end{equation*}

(The interchange of \(\operatorname{Im}\) with the sum is legitimate because the series converges absolutely.)

Next we identify \(g\) in closed form. The power series defining \(g\) has radius of convergence \(1\), so \(g\) is analytic on \(\mathbb{D}\) with \(g(0) = 0\) and

\begin{equation*} g^{\prime}(w) = \sum_{k=0}^{\infty} w^{2k} = \frac{1}{1 - w^2}\qquad (w \in \mathbb{D}). \end{equation*}

Let \(\psi(w) = \dfrac{1+w}{1-w}\) for \(w \in \mathbb{D}\). Then

\begin{equation*} \psi(w) = \frac{(1+w)(1-\overline{w})}{|1-w|^2} = \frac{\bigl(1 - |w|^2\bigr) + 2i\operatorname{Im} w}{|1-w|^2}, \end{equation*}

so \(\operatorname{Re}\psi(w) = \dfrac{1-|w|^2}{|1-w|^2} > 0\) for every \(w \in \mathbb{D}\). Thus \(\psi\) maps \(\mathbb{D}\) into the open right half-plane, where the principal logarithm \(\operatorname{Log}\) is analytic. Hence \(h := \tfrac12 \operatorname{Log}\circ\,\psi\) is analytic on \(\mathbb{D}\), with \(h(0) = \tfrac12\operatorname{Log} 1 = 0\) and, since \(\psi^{\prime}(w) = \frac{2}{(1-w)^2}\),

\begin{equation*} h^{\prime}(w) = \frac{1}{2}\cdot\frac{\psi^{\prime}(w)}{\psi(w)} = \frac{1}{2}\cdot\frac{2}{(1-w)^2}\cdot\frac{1-w}{1+w} = \frac{1}{(1-w)(1+w)} = \frac{1}{1-w^2}. \end{equation*}

So \(g\) and \(h\) are analytic on the connected open set \(\mathbb{D}\), have the same derivative, and agree at \(0\); therefore \(g = h\) on \(\mathbb{D}\).

Taking imaginary parts, and recalling that \(\operatorname{Im}\operatorname{Log}\zeta = \operatorname{Arg}\zeta\) and that \(\operatorname{Arg}(a + bi) = \arctan\frac{b}{a}\) whenever \(a > 0\), we get from the computation of \(\psi(w)\) above (whose real part is \(\frac{1-|w|^2}{|1-w|^2} > 0\) and whose imaginary part is \(\frac{2\operatorname{Im}w}{|1-w|^2}\))

\begin{equation*} \operatorname{Im} g(w) = \tfrac12 \operatorname{Arg}\psi(w) = \tfrac12 \arctan\frac{2\operatorname{Im}w}{1 - |w|^2}. \end{equation*}

Therefore

\begin{equation*} (\mathcal{P}_r f)(z) = \frac{4}{\pi}\cdot\frac{1}{2}\arctan\frac{2\operatorname{Im}w}{1-|w|^2} = \frac{2}{\pi}\arctan\frac{2r\operatorname{Im}z}{1 - r^2}, \end{equation*}

as claimed.

(c) Fix \(z \in \partial\mathbb{D}\) and let \(r \uparrow 1\); note \(1 - r^2 \downarrow 0\) with \(1 - r^2 > 0\), and \(2r \to 2\).

If \(\operatorname{Im}z > 0\), then \(\dfrac{2r\operatorname{Im}z}{1-r^2}\to +\infty\), so since \(\lim_{s\to+\infty}\arctan s = \frac{\pi}{2}\) we get \((\mathcal{P}_r f)(z) \to \frac{2}{\pi}\cdot\frac{\pi}{2} = 1 = f(z)\).

If \(\operatorname{Im}z < 0\), then \(\dfrac{2r\operatorname{Im}z}{1-r^2}\to -\infty\), so \((\mathcal{P}_r f)(z) \to \frac{2}{\pi}\cdot\bigl(-\frac{\pi}{2}\bigr) = -1 = f(z)\).

If \(\operatorname{Im}z = 0\) (that is, \(z = 1\) or \(z = -1\)), then the argument of \(\arctan\) is \(0\) for every \(r\), so \((\mathcal{P}_r f)(z) = 0\) for all \(r \in [0,1)\), and the limit is \(0 = f(z)\).

In all three cases \(\lim_{r\uparrow1}(\mathcal{P}_r f)(z) = f(z)\). (For \(z \neq \pm 1\) this also follows from Exercise 4 in this section, since \(f\) is continuous at such \(z\); for \(z = \pm 1\) it follows from Exercise 5, since \(f\) has a jump discontinuity there with one-sided limits \(1\) and \(-1\), whose average is \(0 = f(\pm 1)\).)

(d) Suppose \(r \in [\sqrt{2}-1, \, 1)\). Then \(s_r := \dfrac{1-r^2}{2r}\) satisfies \(0 < s_r \leq 1\): positivity is clear, and \(s_r \leq 1\) is equivalent to \(1 - r^2 \leq 2r\), i.e. to \(r^2 + 2r - 1 \geq 0\), i.e. to \(r \geq \sqrt{2}-1\). Define

\begin{equation*} z_r = \sqrt{1 - s_r^{\,2}} + i\, s_r \in \partial\mathbb{D}. \end{equation*}

Then \(\operatorname{Im}z_r = s_r > 0\), so \(f(z_r) = 1\), while by (b)

\begin{equation*} (\mathcal{P}_r f)(z_r) = \frac{2}{\pi}\arctan\frac{2r s_r}{1-r^2} = \frac{2}{\pi}\arctan 1 = \frac{2}{\pi}\cdot\frac{\pi}{4} = \frac{1}{2}. \end{equation*}

Hence

\begin{equation*} \|f - \mathcal{P}_r f\|_\infty \geq \bigl|f(z_r) - (\mathcal{P}_r f)(z_r)\bigr| = \frac{1}{2} \end{equation*}

for every \(r \in [\sqrt2 - 1, 1)\). Thus \(\|f - \mathcal{P}_r f\|_\infty\) does not tend to \(0\) as \(r \uparrow 1\), so \(\mathcal{P}_r f\) does not converge uniformly to \(f\) on \(\partial\mathbb{D}\). \(\blacksquare\)

Remark. Part (d) also follows abstractly: as noted after 11.11, each \(\mathcal{P}_r f\) is continuous on \(\partial\mathbb{D}\), and a uniform limit of continuous functions is continuous; but \(f\) is discontinuous at \(1\) (indeed \(f(e^{it}) \to 1\) as \(t \downarrow 0\) while \(f(1) = 0\)), so no uniform convergence is possible. Comparing with 11.18, which asserts uniform convergence \(\mathcal{P}_r f \to f\) for continuous \(f\), this exercise shows that the continuity hypothesis in 11.18 cannot be dropped: pointwise convergence, as verified in (c), is all that survives here.

Exercises 11B

Problem (11B.1)

Show that the family \(\{e_k\}_{k \in \mathbb{Z}}\) of trigonometric functions defined by 11.1 is an orthonormal basis of \(L^2\big((-\pi, \pi]\big)\).

Solution

Recall 11.1: for \(k \in \mathbb{Z}\) the function \(e_k \colon (-\pi, \pi] \to \mathbb{R}\) is defined by

\begin{equation*} e_k(t) = \begin{cases} \frac{1}{\sqrt{\pi}} \sin(kt) & \text{if } k > 0,\\[2pt] \frac{1}{\sqrt{2\pi}} & \text{if } k = 0,\\[2pt] \frac{1}{\sqrt{\pi}} \cos(kt) & \text{if } k < 0. \end{cases} \end{equation*}

Here \(L^2\big((-\pi,\pi]\big)\) carries ordinary Lebesgue measure on \((-\pi, \pi]\), and (as throughout Chapter 11) the scalar field is \(\mathbb{C}\).

Step 1: transfer \(L^2(\partial \mathbb{D})\) to \(L^2\big((-\pi,\pi]\big)\) by a unitary map.

Let \(\varphi \colon (-\pi, \pi] \to \partial \mathbb{D}\) be the bijection \(\varphi(t) = e^{it}\) of 11.2. By the definition 11.4 of the measurable subsets of \(\partial \mathbb{D}\), a set \(E \subseteq \partial\mathbb{D}\) is measurable if and only if \(\varphi^{-1}(E)\) is a Borel subset of \((-\pi, \pi]\). Consequently a function \(f \colon \partial \mathbb{D} \to \mathbb{C}\) is measurable if and only if \(f \circ \varphi\) is Borel measurable on \((-\pi, \pi]\), and \(g \mapsto g \circ \varphi^{-1}\) is a bijection from the Borel measurable functions on \((-\pi, \pi]\) onto the measurable functions on \(\partial \mathbb{D}\).

Define \(U \colon L^2(\partial \mathbb{D}) \to L^2\big((-\pi, \pi]\big)\) by

\begin{equation*} (Uf)(t) = \frac{1}{\sqrt{2\pi}} \, f(e^{it}). \end{equation*}

The transfer of integration recorded just after 11.4 gives, for \(f, h \in L^2(\partial \mathbb{D})\),

\begin{equation*} \langle Uf, Uh \rangle = \int_{-\pi}^{\pi} f(e^{it}) \overline{h(e^{it})} \, \frac{dt}{2\pi} = \int_{\partial \mathbb{D}} f \overline{h} \, d\sigma = \langle f, h \rangle . \end{equation*}

In particular \(\|Uf\|_2 = \|f\|_2\), so \(U\) is a well-defined linear isometry. It is surjective: given \(g \in L^2\big((-\pi,\pi]\big)\), the function \(f = \sqrt{2\pi} \, (g \circ \varphi^{-1})\) is measurable on \(\partial \mathbb{D}\) by the previous paragraph, lies in \(L^2(\partial \mathbb{D})\) by the displayed change of variables, and satisfies \(Uf = g\). Thus \(U\) is a unitary map between Hilbert spaces; in particular \(U\) and \(U^{-1}\) are continuous, so \(U\) carries closures of subspaces onto closures of their images.

By 11.30 the family \(\{z^n\}_{n \in \mathbb{Z}}\) is an orthonormal basis of \(L^2(\partial \mathbb{D})\). Put

\begin{equation*} f_n(t) = U(z^n)(t) = \frac{e^{int}}{\sqrt{2\pi}} \qquad (n \in \mathbb{Z}). \end{equation*}

Because \(U\) preserves inner products, \(\{f_n\}_{n \in \mathbb{Z}}\) is an orthonormal family; because \(U\) is a surjective homeomorphism,

\begin{equation*} \overline{\operatorname{span}} \{f_n\}_{n \in \mathbb{Z}} = U\Big( \overline{\operatorname{span}} \{z^n\}_{n \in \mathbb{Z}} \Big) = U\big(L^2(\partial \mathbb{D})\big) = L^2\big((-\pi,\pi]\big). \end{equation*}

So \(\{f_n\}_{n \in \mathbb{Z}}\) is an orthonormal basis of \(L^2\big((-\pi,\pi]\big)\) in the sense of 8.61.

Step 2: \(\{e_k\}\) and \(\{f_n\}\) have the same span.

First, \(e_0 = \frac{1}{\sqrt{2\pi}} = f_0\). Next fix \(n \in \mathbb{Z}^{+}\). By the definition 11.1 we have \(e_n(t) = \frac{1}{\sqrt{\pi}}\sin(nt)\) and, since \(-n < 0\) and \(\cos(-nt) = \cos(nt)\), also \(e_{-n}(t) = \frac{1}{\sqrt{\pi}} \cos(nt)\). Hence \(\cos(nt) = \sqrt{\pi}\, e_{-n}(t)\) and \(\sin(nt) = \sqrt{\pi} \, e_n(t)\), and therefore

\begin{equation*} f_{n} = \frac{\cos(nt) + i \sin(nt)}{\sqrt{2\pi}} = \frac{e_{-n} + i e_{n}}{\sqrt{2}}, \qquad f_{-n} = \frac{\cos(nt) - i \sin(nt)}{\sqrt{2\pi}} = \frac{e_{-n} - i e_{n}}{\sqrt{2}}. \end{equation*}

Solving these two equations gives

\begin{equation*} e_{-n} = \frac{f_n + f_{-n}}{\sqrt{2}}, \qquad e_{n} = \frac{f_n - f_{-n}}{i \sqrt{2}} . \end{equation*}

Thus each \(e_k\) is a finite linear combination of the \(f_n\) and each \(f_n\) is a finite linear combination of the \(e_k\), so

\begin{equation*} \operatorname{span}\{e_k\}_{k \in \mathbb{Z}} = \operatorname{span}\{f_n\}_{n \in \mathbb{Z}}, \end{equation*}

and hence the two families have the same closure of the span, namely all of \(L^2\big((-\pi,\pi]\big)\).

Step 3: \(\{e_k\}\) is orthonormal.

This now follows from the orthonormality of \(\{f_n\}\) and the formulas above. Indeed \(\|e_0\| = \|f_0\| = 1\), and for \(n \in \mathbb{Z}^{+}\),

\begin{equation*} \|e_{-n}\|^2 = \tfrac{1}{2}\big(\|f_n\|^2 + \|f_{-n}\|^2\big) = 1, \qquad \|e_{n}\|^2 = \tfrac{1}{|i|^2 \cdot 2}\big(\|f_n\|^2 + \|f_{-n}\|^2\big) = 1, \end{equation*}

using \(\langle f_n, f_{-n}\rangle = 0\). Also

\begin{equation*} \langle e_n, e_{-n}\rangle = \Big\langle \frac{f_n - f_{-n}}{i\sqrt{2}}, \frac{f_n + f_{-n}}{\sqrt{2}} \Big\rangle = \frac{1}{2i}\big(\|f_n\|^2 - \|f_{-n}\|^2\big) = 0 . \end{equation*}

If \(j, k \in \mathbb{Z}\) with \(|j| \ne |k|\), then \(e_j\) is a linear combination of \(f_{|j|}, f_{-|j|}\) (or of \(f_0\) when \(j = 0\)) and \(e_k\) is a linear combination of \(f_{|k|}, f_{-|k|}\); these two pairs are disjoint, so \(\langle e_j, e_k \rangle = 0\). Combining the cases shows that \(\langle e_j, e_k\rangle\) equals \(1\) if \(j = k\) and \(0\) otherwise.

By Steps 2 and 3, \(\{e_k\}_{k \in \mathbb{Z}}\) is an orthonormal family whose span is dense in \(L^2\big((-\pi,\pi]\big)\). By the definition 8.61 of an orthonormal basis, \(\{e_k\}_{k \in \mathbb{Z}}\) is an orthonormal basis of \(L^2\big((-\pi,\pi]\big)\).

Finally, a remark for the classical real setting. Every \(e_k\) is real valued, and every real-valued \(g \in L^2\big((-\pi,\pi]\big)\) can be approximated in norm by complex linear combinations \(\sum \alpha_k e_k\); taking real parts and using \(\|\operatorname{Re} h\|_2 \le \|h\|_2\) shows that \(g\) is approximated equally well by \(\sum (\operatorname{Re}\alpha_k) e_k\), a real linear combination. Hence \(\{e_k\}_{k \in \mathbb{Z}}\) is also an orthonormal basis of the real Hilbert space \(L^2\big((-\pi,\pi]\big)\).

Problem (11B.2)

Use the result of Exercise 12(a) in Section 11A to show that

\begin{equation*} 1 + \frac{1}{3^2} + \frac{1}{5^2} + \frac{1}{7^2} + \cdots = \frac{\pi^2}{8}. \end{equation*}

Solution

Let \(f \colon \partial \mathbb{D} \to \mathbb{R}\) be the function of Exercise 12 in Section 11A:

\begin{equation*} f(z) = \begin{cases} 1 & \text{if } \operatorname{Im} z > 0,\\ -1 & \text{if } \operatorname{Im} z < 0,\\ 0 & \text{if } \operatorname{Im} z = 0 . \end{cases} \end{equation*}

By Exercise 12(a) in Section 11A,

\begin{equation*} \widehat{f}(n) = \begin{cases} -\dfrac{2i}{n\pi} & \text{if } n \text{ is odd},\\[6pt] 0 & \text{if } n \text{ is even}. \end{cases} \end{equation*}

(For completeness, here is the computation. With \(f(e^{it}) = 1\) for \(t \in (0,\pi)\), \(f(e^{it}) = -1\) for \(t \in (-\pi, 0)\), and \(f(e^{it}) = 0\) for \(t \in \{0, \pi\}\), we get \(\widehat{f}(0) = \frac{1}{2\pi}(\pi - \pi) = 0\), while for \(n \ne 0\)

\begin{equation*} \widehat{f}(n) = \frac{1}{2\pi}\left( \int_0^{\pi} e^{-int} \, dt - \int_{-\pi}^{0} e^{-int} \, dt \right) = \frac{1}{2\pi} \cdot \frac{2(-1)^n - 2}{-in}, \end{equation*}

which is \(0\) for even \(n\) and equals \(\frac{1}{2\pi} \cdot \frac{-4}{-in} = -\frac{2i}{n\pi}\) for odd \(n\).)

Now compute \(\|f\|_2\). The set \(\{z \in \partial\mathbb{D} : \operatorname{Im} z = 0\} = \{1, -1\}\) consists of two points and hence has \(\sigma\)-measure \(0\) (by 11.4, since its inverse image under \(t \mapsto e^{it}\) is the two-point set \(\{0, \pi\}\), which has Lebesgue measure \(0\)). Thus \(|f(z)| = 1\) for almost every \(z \in \partial \mathbb{D}\), so \(f \in L^2(\partial \mathbb{D})\) and

\begin{equation*} \|f\|_2^2 = \int_{\partial \mathbb{D}} |f|^2 \, d\sigma = \sigma(\partial \mathbb{D}) = 1 . \end{equation*}

By 11.30 the family \(\{z^n\}_{n \in \mathbb{Z}}\) is an orthonormal basis of \(L^2(\partial \mathbb{D})\), and \(\langle f, z^n \rangle = \int_{\partial\mathbb{D}} f(z) \overline{z^n} \, d\sigma(z) = \widehat{f}(n)\) by the definition 11.7 of the Fourier coefficients. Hence Parseval’s identity [8.63(c)] gives

\begin{equation*} \sum_{n = -\infty}^{\infty} |\widehat{f}(n)|^2 = \|f\|_2^2 = 1 . \end{equation*}

On the other hand, only the odd indices contribute, and the terms for \(n\) and \(-n\) are equal, so

\begin{equation*} \sum_{n = -\infty}^{\infty} |\widehat{f}(n)|^2 = \sum_{n \text{ odd}} \frac{4}{n^2 \pi^2} = 2 \sum_{k=0}^{\infty} \frac{4}{(2k+1)^2 \pi^2} = \frac{8}{\pi^2} \sum_{k=0}^{\infty} \frac{1}{(2k+1)^2} . \end{equation*}

(All terms are nonnegative, so the rearrangement into pairs is legitimate; equivalently, the unordered sum of a family of nonnegative numbers is unchanged by regrouping.)

Combining the two displays gives

\begin{equation*} \frac{8}{\pi^2} \sum_{k=0}^{\infty} \frac{1}{(2k+1)^2} = 1, \end{equation*}

and therefore

\begin{equation*} 1 + \frac{1}{3^2} + \frac{1}{5^2} + \frac{1}{7^2} + \cdots = \sum_{k=0}^{\infty} \frac{1}{(2k+1)^2} = \frac{\pi^2}{8}, \end{equation*}

as desired.

Problem (11B.3)

Use techniques similar to Example 11.32 to evaluate \(\displaystyle\sum_{n=1}^{\infty} \frac{1}{n^4}\).

[If you feel industrious, you may also want to evaluate \(\sum_{n=1}^{\infty} 1/n^6\). Similar techniques work to evaluate \(\sum_{n=1}^{\infty} 1/n^k\) for each positive even integer \(k\). You can become famous if you figure out how to evaluate \(\sum_{n=1}^{\infty} 1/n^3\), which currently is an open question.]

Solution

Throughout, \(\zeta(k)\) abbreviates \(\sum_{n=1}^{\infty} 1/n^k\); Example 11.32 gives \(\zeta(2) = \pi^2/6\).

Evaluation of \(\zeta(4)\). Define \(f \in L^2(\partial \mathbb{D})\) by

\begin{equation*} f(e^{it}) = t^2 \qquad \text{for } t \in (-\pi, \pi] . \end{equation*}

The function \(t \mapsto t^2\) is continuous, hence Borel measurable, on \((-\pi,\pi]\), so \(f\) is a measurable function on \(\partial \mathbb{D}\) by 11.4; it is bounded and \(\sigma(\partial \mathbb{D}) = 1\), so indeed \(f \in L^2(\partial \mathbb{D})\).

First,

\begin{equation*} \widehat{f}(0) = \int_{-\pi}^{\pi} t^2 \, \frac{dt}{2\pi} = \frac{1}{2\pi} \cdot \frac{2\pi^3}{3} = \frac{\pi^2}{3} . \end{equation*}

Now suppose \(n \in \mathbb{Z} \setminus \{0\}\). Because \(t \mapsto t^2 \sin(nt)\) is odd and \(t \mapsto t^2\cos(nt)\) is even,

\begin{equation*} \widehat{f}(n) = \int_{-\pi}^{\pi} t^2 e^{-int} \, \frac{dt}{2\pi} = \frac{1}{2\pi} \int_{-\pi}^{\pi} t^2 \cos(nt) \, dt = \frac{1}{\pi} \int_{0}^{\pi} t^2 \cos(nt) \, dt . \end{equation*}

Two integrations by parts (or direct differentiation of the antiderivative) give

\begin{equation*} \int t^2 \cos(nt) \, dt = \frac{t^2 \sin(nt)}{n} + \frac{2t\cos(nt)}{n^2} - \frac{2 \sin(nt)}{n^3}, \end{equation*}

which vanishes at \(t = 0\) and equals \(2\pi(-1)^n/n^2\) at \(t = \pi\) (since \(\sin(n\pi) = 0\) and \(\cos(n\pi) = (-1)^n\)). Hence

\begin{equation*} \widehat{f}(n) = \frac{1}{\pi} \cdot \frac{2\pi (-1)^n}{n^2} = \frac{2(-1)^n}{n^2} \qquad \text{for } n \ne 0 . \end{equation*}

Therefore

\begin{equation*} \sum_{n=-\infty}^{\infty} |\widehat{f}(n)|^2 = \frac{\pi^4}{9} + 2\sum_{n=1}^{\infty} \frac{4}{n^4} = \frac{\pi^4}{9} + 8\,\zeta(4) . \end{equation*}

On the other hand,

\begin{equation*} \|f\|_2^2 = \int_{-\pi}^{\pi} t^4 \, \frac{dt}{2\pi} = \frac{1}{2\pi} \cdot \frac{2\pi^5}{5} = \frac{\pi^4}{5} . \end{equation*}

Because \(\{z^n\}_{n \in \mathbb{Z}}\) is an orthonormal basis of \(L^2(\partial \mathbb{D})\) (by 11.30) and \(\langle f, z^n \rangle = \widehat{f}(n)\), Parseval’s identity [8.63(c)] equates the two displays above:

\begin{equation*} \frac{\pi^4}{9} + 8\,\zeta(4) = \frac{\pi^4}{5} . \end{equation*}

Thus \(8\,\zeta(4) = \pi^4\big(\tfrac15 - \tfrac19\big) = \tfrac{4\pi^4}{45}\), and so

\begin{equation*} \sum_{n=1}^{\infty} \frac{1}{n^4} = \frac{\pi^4}{90} . \end{equation*}

Evaluation of \(\zeta(6)\). Now define \(g \in L^2(\partial \mathbb{D})\) by \(g(e^{it}) = t^3\) for \(t \in (-\pi,\pi]\); as above \(g\) is measurable and bounded, hence in \(L^2(\partial \mathbb{D})\).

Since \(t \mapsto t^3\) is odd, \(\widehat{g}(0) = 0\). For \(n \in \mathbb{Z}\setminus\{0\}\), oddness of \(t^3 \cos(nt)\) and evenness of \(t^3 \sin(nt)\) give

\begin{equation*} \widehat{g}(n) = \frac{1}{2\pi}\int_{-\pi}^{\pi} t^3 e^{-int} \, dt = \frac{-i}{2\pi} \int_{-\pi}^{\pi} t^3 \sin(nt) \, dt = \frac{-i}{\pi} \int_{0}^{\pi} t^3 \sin(nt) \, dt . \end{equation*}

Repeated integration by parts yields the antiderivative

\begin{equation*} \int t^3 \sin(nt) \, dt = -\frac{t^3 \cos(nt)}{n} + \frac{3t^2 \sin(nt)}{n^2} + \frac{6t\cos(nt)}{n^3} - \frac{6\sin(nt)}{n^4}, \end{equation*}

as one checks by differentiating. It vanishes at \(t = 0\), and at \(t = \pi\) it equals \((-1)^n\big(\frac{6\pi}{n^3} - \frac{\pi^3}{n}\big)\). Hence

\begin{equation*} \widehat{g}(n) = \frac{-i}{\pi} \cdot (-1)^n\left(\frac{6\pi}{n^3} - \frac{\pi^3}{n}\right) = i(-1)^n\left(\frac{\pi^2}{n} - \frac{6}{n^3}\right), \end{equation*}

so that

\begin{equation*} |\widehat{g}(n)|^2 = \frac{\pi^4}{n^2} - \frac{12\pi^2}{n^4} + \frac{36}{n^6} . \end{equation*}

Summing over \(n \ne 0\) (each term appears twice, for \(n\) and \(-n\)) and using \(\zeta(2) = \pi^2/6\) and \(\zeta(4) = \pi^4/90\):

\begin{equation*} \begin{aligned} \sum_{n=-\infty}^{\infty} |\widehat{g}(n)|^2 &= 2\big(\pi^4 \zeta(2) - 12\pi^2 \zeta(4) + 36\, \zeta(6)\big)\\ &= 2\left(\frac{\pi^6}{6} - \frac{12\pi^6}{90} + 36\,\zeta(6)\right) = 2\left(\frac{\pi^6}{30} + 36\,\zeta(6)\right) = \frac{\pi^6}{15} + 72\,\zeta(6) . \end{aligned} \end{equation*}

(The three separate series converge absolutely, so splitting the sum is legitimate.) Meanwhile

\begin{equation*} \|g\|_2^2 = \int_{-\pi}^{\pi} t^6 \, \frac{dt}{2\pi} = \frac{1}{2\pi}\cdot \frac{2\pi^7}{7} = \frac{\pi^6}{7} . \end{equation*}

Parseval’s identity [8.63(c)] again equates these:

\begin{equation*} \frac{\pi^6}{15} + 72\, \zeta(6) = \frac{\pi^6}{7}, \end{equation*}

so \(72\,\zeta(6) = \pi^6\big(\tfrac17 - \tfrac1{15}\big) = \tfrac{8\pi^6}{105}\), giving

\begin{equation*} \sum_{n=1}^{\infty} \frac{1}{n^6} = \frac{\pi^6}{945} . \end{equation*}

The same method evaluates \(\zeta(k)\) for every even positive integer \(k\). Fix \(m \in \mathbb{Z}^{+}\) and let \(f_m(e^{it}) = t^{m}\) for \(t \in (-\pi,\pi]\); as above \(f_m\) is bounded and measurable, hence lies in \(L^2(\partial\mathbb{D})\). Write \(I_m = \int_{-\pi}^{\pi} t^m e^{-int} \, dt\) for \(n \ne 0\). Integration by parts gives the recursion

\begin{equation*} I_m = (-1)^n \, \frac{\pi^m \big(1 - (-1)^m\big)}{-in} + \frac{m}{in} \, I_{m-1}, \qquad I_0 = 0 , \end{equation*}

using \(e^{-in\pi} = e^{in\pi} = (-1)^n\). An easy induction on \(m\) now shows that

\begin{equation*} \widehat{f_m}(n) = \frac{I_m}{2\pi} = (-1)^n \, i^{\,m} \sum_{\substack{1 \le j \le m \\ j \equiv m \ (\mathrm{mod}\ 2)}} \frac{a_j}{n^{\,j}} \qquad (n \ne 0), \end{equation*}

where the \(a_j\) are real numbers depending only on \(m\). (The cases \(m = 2\) and \(m = 3\) computed above illustrate this.) Every index \(j\) appearing here has the same parity as \(m\), so \(j + k\) is even for any two such indices, and therefore

\begin{equation*} |\widehat{f_m}(n)|^2 = \sum_{j,k} \frac{a_j a_k}{n^{\,j+k}} \end{equation*}

is a finite real combination of \(1/n^2, 1/n^4, \dots, 1/n^{2m}\) in which the coefficient of \(1/n^{2m}\) is \(a_m^2\), and \(a_m \ne 0\) because the recursion multiplies the top-order term by \(m/(in)\) at each step. Summing over \(n \ne 0\) (each of the finitely many series converges absolutely) and equating with

\begin{equation*} \|f_m\|_2^2 = \int_{-\pi}^{\pi} t^{2m} \, \frac{dt}{2\pi} = \frac{\pi^{2m}}{2m+1} \end{equation*}

via Parseval’s identity [8.63(c)] therefore determines \(\zeta(2m)\) from \(\pi^{2m}\) and \(\zeta(2), \zeta(4), \dots, \zeta(2m-2)\). Induction on \(m\) then evaluates every \(\zeta(2m)\), always as a rational multiple of \(\pi^{2m}\).

No such recursion is available for odd exponents, because \(j + k\) above is always even and so the series \(\sum 1/n^3, \sum 1/n^5, \dots\) never appear. No closed form for \(\zeta(3)\) is known.

Problem (11B.4)

Suppose \(f, g \colon \partial \mathbb{D} \to \mathbb{C}\) are measurable functions. Prove that the function \((w, z) \mapsto f(w)\, g(z \overline{w})\) is a measurable function from \(\partial \mathbb{D} \times \partial \mathbb{D}\) to \(\mathbb{C}\).

[Here the \(\sigma\)-algebra on \(\partial \mathbb{D} \times \partial \mathbb{D}\) is the usual product \(\sigma\)-algebra as defined in 5.2.]

Solution

Write \(\mathcal{S}\) for the \(\sigma\)-algebra of measurable subsets of \(\partial \mathbb{D}\) (see 11.4), \(\mathcal{B}\) for the \(\sigma\)-algebra of Borel subsets of \(\mathbb{R}\) that are contained in \((-\pi, \pi]\), and \(\varphi \colon (-\pi, \pi] \to \partial \mathbb{D}\) for the bijection \(\varphi(t) = e^{it}\). By 11.4,

\begin{equation*} \mathcal{S} = \{ E \subseteq \partial \mathbb{D} : \varphi^{-1}(E) \in \mathcal{B}\} = \{\varphi(B) : B \in \mathcal{B}\}, \end{equation*}

the second equality because \(\varphi\) is a bijection. A complex-valued function is measurable exactly when its real and imaginary parts are measurable, so it suffices throughout to work with real-valued functions.

Let \(m \colon \partial \mathbb{D} \times \partial \mathbb{D} \to \partial \mathbb{D}\) be defined by \(m(w,z) = z \overline{w}\). Note that \(m\) does map into \(\partial\mathbb{D}\), since \(|z\overline{w}| = 1\). The heart of the proof is the following claim.

Claim: \(m\) is measurable, in the sense that \(m^{-1}(E) \in \mathcal{S} \otimes \mathcal{S}\) for every \(E \in \mathcal{S}\).

To prove the claim, define \(\psi \colon (-\pi,\pi] \times (-\pi,\pi] \to (-\pi,\pi]\) by letting \(\psi(s,t)\) be the unique element of \((-\pi,\pi]\) congruent to \(t - s\) modulo \(2\pi\). Because \(s, t \in (-\pi, \pi]\) we have \(t - s \in (-2\pi, 2\pi)\), and so explicitly

\begin{equation*} \psi(s,t) = \begin{cases} t - s + 2\pi & \text{if } t - s \le -\pi,\\ t - s & \text{if } -\pi < t - s \le \pi,\\ t - s - 2\pi & \text{if } t - s > \pi . \end{cases} \end{equation*}

The point of \(\psi\) is the identity

\begin{equation*} \varphi(\psi(s,t)) = e^{i(t-s)} = e^{it} \, \overline{e^{is}} = m(\varphi(s), \varphi(t)), \end{equation*}

valid because \(e^{i\theta}\) has period \(2\pi\).

Sub-step (i): \(\psi^{-1}(B) \in \mathcal{B} \otimes \mathcal{B}\) for every \(B \in \mathcal{B}\). The coordinate functions \((s,t) \mapsto s\) and \((s,t) \mapsto t\) are \(\mathcal{B}\otimes\mathcal{B}\)-measurable, since the inverse image of a Borel set \(A \subseteq \mathbb{R}\) under the first is the measurable rectangle \(\big(A \cap (-\pi,\pi]\big) \times (-\pi,\pi]\), and similarly for the second. Hence for each constant \(c \in \mathbb{R}\) the function \(h_c(s,t) = t - s + c\) is \(\mathcal{B}\otimes\mathcal{B}\)-measurable by 2.46(a), and therefore \(h_c^{-1}(A) \in \mathcal{B}\otimes\mathcal{B}\) for every Borel \(A \subseteq \mathbb{R}\). In particular the three sets

\begin{equation*} D_1 = h_0^{-1}\big((-\infty, -\pi]\big), \quad D_2 = h_0^{-1}\big((-\pi, \pi]\big), \quad D_3 = h_0^{-1}\big((\pi, \infty)\big) \end{equation*}

belong to \(\mathcal{B}\otimes\mathcal{B}\), and for \(B \in \mathcal{B}\),

\begin{equation*} \psi^{-1}(B) = \big(D_1 \cap h_{2\pi}^{-1}(B)\big) \cup \big(D_2 \cap h_0^{-1}(B)\big) \cup \big(D_3 \cap h_{-2\pi}^{-1}(B)\big) \in \mathcal{B}\otimes\mathcal{B} . \end{equation*}

Sub-step (ii): if \(C \in \mathcal{B}\otimes\mathcal{B}\), then \((\varphi \times \varphi)( C) \in \mathcal{S}\otimes\mathcal{S}\). Let

\begin{equation*} \mathcal{C} = \big\{ C \subseteq (-\pi,\pi] \times (-\pi,\pi] : (\varphi\times\varphi)( C) \in \mathcal{S}\otimes\mathcal{S} \big\}. \end{equation*}

Because \(\varphi \times \varphi\) is a bijection onto \(\partial\mathbb{D}\times\partial\mathbb{D}\), taking images under it commutes with complements and with arbitrary unions; hence \(\mathcal{C}\) is a \(\sigma\)-algebra. If \(B_1, B_2 \in \mathcal{B}\), then \((\varphi\times\varphi)(B_1 \times B_2) = \varphi(B_1) \times \varphi(B_2)\), which is a measurable rectangle in \(\mathcal{S}\otimes\mathcal{S}\) because \(\varphi(B_1), \varphi(B_2) \in \mathcal{S}\). So \(\mathcal{C}\) is a \(\sigma\)-algebra containing every measurable rectangle of \(\mathcal{B}\otimes\mathcal{B}\), and therefore \(\mathcal{B}\otimes\mathcal{B} \subseteq \mathcal{C}\), which is the assertion of sub-step (ii).

Now let \(E \in \mathcal{S}\), so \(B := \varphi^{-1}(E) \in \mathcal{B}\). Using that \(\varphi\times\varphi\) is a bijection and the identity \(\varphi \circ \psi = m \circ (\varphi\times\varphi)\), we get

\begin{equation*} m^{-1}(E) = (\varphi\times\varphi)\big( (\varphi\times\varphi)^{-1}(m^{-1}(E)) \big) = (\varphi\times\varphi)\big( \psi^{-1}(\varphi^{-1}(E)) \big) = (\varphi\times\varphi)\big(\psi^{-1}(B)\big). \end{equation*}

By sub-step (i), \(\psi^{-1}(B) \in \mathcal{B}\otimes\mathcal{B}\); by sub-step (ii), its image under \(\varphi\times\varphi\) lies in \(\mathcal{S}\otimes\mathcal{S}\). This proves the claim.

With the claim in hand the rest is routine.

First, \((w,z) \mapsto f(w)\) is \(\mathcal{S}\otimes\mathcal{S}\)-measurable: for \(a \in \mathbb{R}\) the inverse image of \((a,\infty)\) under \((w,z) \mapsto \operatorname{Re} f(w)\) is \((\operatorname{Re} f)^{-1}\big((a,\infty)\big) \times \partial\mathbb{D}\), a measurable rectangle, and the same works for \(\operatorname{Im} f\).

Second, \((w,z) \mapsto g(z\overline{w}) = (g \circ m)(w,z)\) is \(\mathcal{S}\otimes\mathcal{S}\)-measurable: for \(a \in \mathbb{R}\),

\begin{equation*} \big( (\operatorname{Re} g) \circ m \big)^{-1}\big((a,\infty)\big) = m^{-1}\Big( (\operatorname{Re} g)^{-1}\big((a,\infty)\big) \Big) \in \mathcal{S}\otimes\mathcal{S}, \end{equation*}

since \((\operatorname{Re} g)^{-1}\big((a,\infty)\big) \in \mathcal{S}\) by measurability of \(g\) and since \(m\) is measurable by the claim; again the same works for \(\operatorname{Im} g\).

Finally, write \(f(w) = u_1(w,z) + i v_1(w,z)\) and \(g(z\overline{w}) = u_2(w,z) + i v_2(w,z)\), where \(u_1, v_1, u_2, v_2\) are the real-valued \(\mathcal{S}\otimes\mathcal{S}\)-measurable functions just produced. Then

\begin{equation*} f(w) g(z\overline{w}) = (u_1 u_2 - v_1 v_2) + i (u_1 v_2 + v_1 u_2), \end{equation*}

and both the real and imaginary parts on the right are \(\mathcal{S}\otimes\mathcal{S}\)-measurable by 2.46(a) (products and sums of measurable real-valued functions are measurable). Hence \((w,z) \mapsto f(w) g(z\overline{w})\) is a measurable function from \(\partial\mathbb{D}\times\partial\mathbb{D}\) to \(\mathbb{C}\), as desired.

Problem (11B.5)

Where does the proof of 11.42 fail when \(p = \infty\)?

Solution

The proof of 11.42 fails at exactly one place: its very first step, the choice of a continuous function \(g\) on \(\partial \mathbb{D}\) with \(\|f - g\|_p < \varepsilon\). Continuous functions are dense in \(L^p(\partial \mathbb{D})\) for \(1 \le p < \infty\), but they are not dense in \(L^{\infty}(\partial \mathbb{D})\).

Every other step of the proof is valid for \(p = \infty\):

  • the triangle inequality \(\|f - P_r f\|_{\infty} \le \|f - g\|_{\infty} + \|g - P_r g\|_{\infty} + \|P_r g - P_r f\|_{\infty}\) holds in any normed space;
  • 11.18 supplies \(R \in [0,1)\) with \(\|g - P_r g\|_{\infty} < \varepsilon\) for \(r \in (R,1)\) whenever \(g\) is continuous on \(\partial\mathbb{D}\), with no reference to \(p\);
  • \(P_r(g - f) = P_r * (g-f) = (g-f) * P_r\), by 11.35 and the commutativity of convolution (11.41), which requires only that both functions lie in \(L^1(\partial\mathbb{D})\);
  • 11.38 is stated and proved for all \(p \in [1, \infty]\), so \(\|P_r * (g-f)\|_{\infty} \le \|P_r\|_1 \|g - f\|_{\infty}\);
  • \(\|P_r\|_1 = 1\), by 11.16(a) and 11.16(b).

So the whole argument goes through for \(p = \infty\) once one has a continuous \(g\) close to \(f\) in \(L^{\infty}\) norm, and it is precisely the existence of such a \(g\) that fails.

That the failure is genuine, and not an artifact of the proof. Let \(f\) be the function of Exercise 12 in Section 11A, so \(f(z) = 1\) if \(\operatorname{Im} z > 0\), \(f(z) = -1\) if \(\operatorname{Im} z < 0\), and \(f(z) = 0\) if \(\operatorname{Im} z = 0\). We show that \(\|f - g\|_{\infty} \ge \frac{1}{2}\) for every continuous \(g \colon \partial\mathbb{D} \to \mathbb{C}\).

First note that every nonempty relatively open subset \(V\) of \(\partial\mathbb{D}\) has \(\sigma(V) > 0\): the set \(\{t \in (-\pi,\pi] : e^{it} \in V\}\) is a nonempty relatively open subset of \((-\pi,\pi]\) and hence contains a nondegenerate interval, so it has positive Lebesgue measure, and now use 11.4.

Suppose \(g\) is continuous on \(\partial\mathbb{D}\) and \(\|f - g\|_{\infty} < \frac{1}{2}\). Let \(N = \{z \in \partial\mathbb{D} : |f(z) - g(z)| \ge \frac12\}\); then \(\sigma(N) = 0\). Let \(A_{+} = \{z \in \partial\mathbb{D} : \operatorname{Im} z > 0\}\) and let \(z_0 \in \partial\mathbb{D}\) satisfy \(\operatorname{Im} z_0 \ge 0\). Every relatively open neighborhood \(V\) of \(z_0\) meets \(A_{+}\) in a nonempty relatively open set, which therefore has positive measure and so is not contained in \(N\). Hence there are points \(z \in A_{+} \setminus N\) arbitrarily close to \(z_0\), and for each such \(z\) we have \(|g(z) - 1| = |g(z) - f(z)| < \frac12\). Letting \(z \to z_0\) and using the continuity of \(g\) gives \(|g(z_0) - 1| \le \frac12\). The symmetric argument with \(A_{-} = \{z : \operatorname{Im} z < 0\}\) gives \(|g(z_0) + 1| \le \frac12\) for every \(z_0 \in \partial\mathbb{D}\) with \(\operatorname{Im} z_0 \le 0\). Taking \(z_0 = 1\), which satisfies both hypotheses, yields

\begin{equation*} 2 = |(g(1) + 1) - (g(1) - 1)| \le |g(1)+1| + |g(1)-1| \le \tfrac12 + \tfrac12 = 1, \end{equation*}

a contradiction. Hence no continuous \(g\) is within \(\frac12\) of \(f\) in \(L^{\infty}(\partial\mathbb{D})\).

Moreover, the conclusion of 11.42 itself is false for \(p = \infty\), as the text notes by referring to Exercise 12(d) in Section 11A. Indeed, by Exercise 12(b),

\begin{equation*} (P_r f)(e^{it}) = \frac{2}{\pi} \arctan \frac{2 r \sin t}{1 - r^2} . \end{equation*}

Fix \(r \in [\frac12, 1)\) and let \(t \in (0, 1-r^2]\). Then \(0 < \frac{2r\sin t}{1-r^2} \le \frac{2t}{1-r^2} \le 2\), so \((P_r f)(e^{it}) \le \frac{2}{\pi}\arctan 2\), while \(f(e^{it}) = 1\). Thus

\begin{equation*} |f - P_r f| \ge 1 - \frac{2}{\pi}\arctan 2 > 0 \end{equation*}

on the arc \(\{e^{it} : 0 < t \le 1 - r^2\}\), which has positive \(\sigma\)-measure. Hence \(\|f - P_r f\|_{\infty} \ge 1 - \frac{2}{\pi}\arctan 2\) for all \(r \in [\frac12, 1)\), and \(\|f - P_r f\|_{\infty}\) does not tend to \(0\) as \(r \uparrow 1\).

Problem (11B.6)

Suppose \(f \in L^1(\partial \mathbb{D})\). Prove that \(f\) is real valued (almost everywhere) if and only if \(\widehat{f}(-n) = \overline{\widehat{f}(n)}\) for every \(n \in \mathbb{Z}\).

Solution

Since \(|\overline{f}| = |f|\), the function \(\overline{f}\) also belongs to \(L^1(\partial\mathbb{D})\), so its Fourier coefficients are defined. We first record the identity of Exercise 1 in Section 11A:

\begin{equation*} \widehat{\overline{f}}(n) = \overline{\widehat{f}(-n)} \qquad \text{for every } n \in \mathbb{Z}. \end{equation*}

Indeed, using the definition 11.7 and \(\overline{z^{\,n}} = z^{-n}\) for \(z \in \partial\mathbb{D}\) (valid because \(\overline{z} = 1/z\) when \(|z| = 1\)),

\begin{equation*} \widehat{\overline{f}}(n) = \int_{\partial\mathbb{D}} \overline{f(z)} \, \overline{z^n} \, d\sigma(z) = \overline{\int_{\partial\mathbb{D}} f(z) \, z^{n} \, d\sigma(z)} = \overline{\int_{\partial\mathbb{D}} f(z) \, \overline{z^{-n}} \, d\sigma(z)} = \overline{\widehat{f}(-n)} . \end{equation*}

Suppose first that \(f\) is real valued almost everywhere. Then \(\overline{f} = f\) as elements of \(L^1(\partial\mathbb{D})\), so their Fourier coefficients agree. By the identity above, for every \(n \in \mathbb{Z}\),

\begin{equation*} \widehat{f}(n) = \widehat{\overline{f}}(n) = \overline{\widehat{f}(-n)} , \end{equation*}

and taking complex conjugates gives \(\widehat{f}(-n) = \overline{\widehat{f}(n)}\), as desired.

Conversely, suppose \(\widehat{f}(-n) = \overline{\widehat{f}(n)}\) for every \(n \in \mathbb{Z}\). Taking complex conjugates gives \(\overline{\widehat{f}(-n)} = \widehat{f}(n)\), and hence, by the identity above,

\begin{equation*} \widehat{\overline{f}}(n) = \overline{\widehat{f}(-n)} = \widehat{f}(n) \end{equation*}

for every \(n \in \mathbb{Z}\). Thus the function \(\overline{f} - f\) lies in \(L^1(\partial\mathbb{D})\) and, by the linearity of \(h \mapsto \widehat{h}(n)\), satisfies \(\widehat{(\overline{f} - f)}(n) = 0\) for every \(n \in \mathbb{Z}\). By 11.43, \(\overline{f} - f = 0\), meaning \(\overline{f(z)} = f(z)\) for almost every \(z \in \partial\mathbb{D}\). Hence \(f\) is real valued almost everywhere.

Problem (11B.7)

Suppose \(f \in L^1(\partial \mathbb{D})\). Show that \(f \in L^2(\partial \mathbb{D})\) if and only if \(\displaystyle\sum_{n=-\infty}^{\infty} |\widehat{f}(n)|^2 < \infty\).

Solution

Throughout, recall that for \(h \in L^2(\partial\mathbb{D})\) the Fourier coefficient \(\widehat{h}(n)\) defined in 11.7 coincides with the inner product \(\langle h, z^n\rangle\) in \(L^2(\partial\mathbb{D})\), since \(\langle h, z^n \rangle = \int_{\partial\mathbb{D}} h(z)\overline{z^n} \, d\sigma(z) = \widehat{h}(n)\).

Also recall that \(L^2(\partial\mathbb{D}) \subseteq L^1(\partial\mathbb{D})\): if \(h \in L^2(\partial\mathbb{D})\), then Hölder’s inequality (7.9) with exponents \(2\) and \(2\) gives

\begin{equation*} \|h\|_1 = \int_{\partial\mathbb{D}} |h| \cdot 1 \, d\sigma \le \|h\|_2 \, \|1\|_2 = \|h\|_2 \, \sigma(\partial\mathbb{D})^{1/2} = \|h\|_2 < \infty, \end{equation*}

using \(\sigma(\partial\mathbb{D}) = 1\) from 11.4.

One direction is immediate. Suppose \(f \in L^2(\partial\mathbb{D})\). By 11.30 the family \(\{z^n\}_{n\in\mathbb{Z}}\) is an orthonormal basis of \(L^2(\partial\mathbb{D})\), so Parseval’s identity [8.63(c)] gives

\begin{equation*} \sum_{n=-\infty}^{\infty} |\widehat{f}(n)|^2 = \sum_{n\in\mathbb{Z}} |\langle f, z^n\rangle|^2 = \|f\|_2^2 < \infty . \end{equation*}

(Bessel’s inequality 8.57 alone would already give finiteness here.)

Now the substantive direction. Suppose \(f \in L^1(\partial\mathbb{D})\) and \(\sum_{n=-\infty}^{\infty} |\widehat{f}(n)|^2 < \infty\). The space \(L^2(\partial\mathbb{D})\) is a Hilbert space: it is complete by 7.24, and its norm arises from the inner product \(\langle h, k \rangle = \int_{\partial\mathbb{D}} h \overline{k} \, d\sigma\). Since \(\{z^n\}_{n\in\mathbb{Z}}\) is an orthonormal family in it (11.6), part (a) of 8.58 applies: since \(\{\widehat{f}(n)\}_{n\in\mathbb{Z}}\) is a family of scalars with \(\sum_{n\in\mathbb{Z}} |\widehat{f}(n)|^2 < \infty\), the unordered sum

\begin{equation*} g := \sum_{n \in \mathbb{Z}} \widehat{f}(n) \, z^n \end{equation*}

converges in the norm of \(L^2(\partial\mathbb{D})\) to an element \(g \in L^2(\partial\mathbb{D})\).

Next we check that \(g\) has the same Fourier coefficients as \(f\). Fix \(m \in \mathbb{Z}\). For a finite subset \(\Omega \subseteq \mathbb{Z}\) with \(m \in \Omega\), put \(g_{\Omega} = \sum_{n\in\Omega} \widehat{f}(n) z^n\); orthonormality of \(\{z^n\}\) gives \(\langle g_{\Omega}, z^m \rangle = \widehat{f}(m)\). By the definition of convergence of an unordered sum, there are such finite sets \(\Omega\) with \(\|g - g_{\Omega}\|_2\) arbitrarily small, and by the Cauchy–Schwarz inequality

\begin{equation*} \big| \langle g, z^m \rangle - \widehat{f}(m) \big| = \big| \langle g - g_{\Omega}, z^m \rangle \big| \le \|g - g_{\Omega}\|_2 \, \|z^m\|_2 = \|g - g_{\Omega}\|_2 . \end{equation*}

Letting \(\|g - g_\Omega\|_2 \to 0\) gives \(\widehat{g}(m) = \langle g, z^m\rangle = \widehat{f}(m)\) for every \(m \in \mathbb{Z}\).

Finally, \(g \in L^2(\partial\mathbb{D}) \subseteq L^1(\partial\mathbb{D})\) by the inclusion noted at the start, so \(f - g \in L^1(\partial\mathbb{D})\) and

\begin{equation*} \widehat{(f - g)}(n) = \widehat{f}(n) - \widehat{g}(n) = 0 \end{equation*}

for every \(n \in \mathbb{Z}\). By 11.43 (functions in \(L^1(\partial\mathbb{D})\) are determined by their Fourier coefficients), \(f - g = 0\), that is, \(f = g\) almost everywhere on \(\partial\mathbb{D}\). Since \(g \in L^2(\partial\mathbb{D})\), we conclude \(f \in L^2(\partial\mathbb{D})\).

This completes both directions. Note also that the argument shows \(\|f\|_2^2 = \sum_{n=-\infty}^{\infty} |\widehat{f}(n)|^2\) whenever either side is finite.

Problem (11B.8)

Suppose \(f \in L^2(\partial D)\). Prove that \(|f(z)| = 1\) for almost every \(z \in \partial D\) if and only if

\begin{equation*} \sum_{k=-\infty}^{\infty} \hat{f}(k)\, \overline{\hat{f}(k-n)} = \begin{cases} 1 & \text{if } n = 0,\\ 0 & \text{if } n \neq 0\end{cases} \end{equation*}

for all \(n \in \mathbb{Z}\).

Solution

Throughout, \(\sigma\) is the rotation-invariant probability measure on \(\partial D\) and \(\langle g, h \rangle = \int_{\partial D} g \bar{h} \, d\sigma\) is the inner product of \(L^2(\partial D)\), so that \(\hat{g}(k) = \langle g, z^k \rangle\) for \(g \in L^2(\partial D)\).

Step 1: the series converges absolutely, and it computes a Fourier coefficient of \(|f|^2\).

Since \(f \in L^2(\partial D)\) and \(\sigma\) is a finite measure, the Cauchy–Schwarz inequality gives \(|f|^2 = f \bar{f} \in L^1(\partial D)\). Fix \(n \in \mathbb{Z}\) and set \(g = z^n f\), meaning \(g(z) = z^n f(z)\). Because \(|z^n| = 1\) on \(\partial D\), we have \(g \in L^2(\partial D)\) and, for each \(k \in \mathbb{Z}\),

\begin{equation*} \hat{g}(k) = \int_{\partial D} z^n f(z)\, \overline{z}^{\,k}\, d\sigma(z) = \int_{\partial D} f(z)\, \overline{z}^{\,k-n}\, d\sigma(z) = \hat{f}(k-n), \end{equation*}

where we used \(z^n \overline{z}^{\,k} = \overline{z}^{\,k-n}\), valid since \(z \overline{z} = 1\) on \(\partial D\).

By 11.30 the family \(\{z^k\}_{k \in \mathbb{Z}}\) is an orthonormal basis of \(L^2(\partial D)\), so Parseval’s identity in the form 8.63(b), applied to the pair \(f, g\), gives

\begin{equation*} \langle f, g \rangle = \sum_{k=-\infty}^{\infty} \langle f, z^k \rangle\, \overline{\langle g, z^k \rangle} = \sum_{k=-\infty}^{\infty} \hat{f}(k)\, \overline{\hat{g}(k)} = \sum_{k=-\infty}^{\infty} \hat{f}(k)\, \overline{\hat{f}(k-n)}. \end{equation*}

Bessel’s inequality (8.57) gives \(\sum_k |\hat{f}(k)|^2 \le \|f\|_2^2 < \infty\), and the same holds for \(\hat{g}\); hence by the Cauchy–Schwarz inequality in \(\ell^2(\mathbb{Z})\) the series above converges absolutely. On the other hand,

\begin{equation*} \langle f, g \rangle = \int_{\partial D} f(z)\, \overline{z^n f(z)}\, d\sigma(z) = \int_{\partial D} |f(z)|^2\, \overline{z}^{\,n}\, d\sigma(z) = \widehat{|f|^2}(n). \end{equation*}

Thus for every \(n \in \mathbb{Z}\),

\begin{equation*} \widehat{|f|^2}(n) = \sum_{k=-\infty}^{\infty} \hat{f}(k)\, \overline{\hat{f}(k-n)}. \end{equation*}

Step 2: the equivalence.

Let \(1\) denote the constant function \(1\) on \(\partial D\), which is \(z^0\). Since \(\{z^n\}_{n \in \mathbb{Z}}\) is an orthonormal family in \(L^2(\partial D)\) (11.6),

\begin{equation*} \hat{1}(n) = \int_{\partial D} \overline{z}^{\,n}\, d\sigma(z) = \langle z^0, z^n \rangle = \begin{cases} 1 & \text{if } n = 0,\\ 0 & \text{if } n \neq 0.\end{cases} \end{equation*}

Suppose first that \(|f(z)| = 1\) for almost every \(z \in \partial D\). Then \(|f|^2 = 1\) almost everywhere, so \(\widehat{|f|^2}(n) = \hat{1}(n)\) for every \(n\), which by Step 1 is exactly the displayed condition of the exercise.

Conversely, suppose the displayed condition holds for all \(n \in \mathbb{Z}\). By Step 1 this says \(\widehat{|f|^2}(n) = \hat{1}(n)\) for every \(n \in \mathbb{Z}\). The function \(|f|^2 - 1\) lies in \(L^1(\partial D)\) and, by 11.9(a) and 11.9(b), all of its Fourier coefficients vanish. Hence \(|f|^2 - 1 = 0\) by 11.43; that is, \(|f(z)|^2 = 1\), and therefore \(|f(z)| = 1\), for almost every \(z \in \partial D\).

Problem (11B.9)

For this exercise, for each \(r \in [0,1)\) think of \(\mathcal{P}_r\) as an operator on \(L^2(\partial D)\).

  • (a) Show that \(\mathcal{P}_r\) is a self-adjoint compact operator for each \(r \in [0,1)\).
  • (b) For each \(r \in [0,1)\), find all eigenvalues and eigenvectors of \(\mathcal{P}_r\).
  • (c) Prove or disprove: \(\lim_{r \uparrow 1} \|I - \mathcal{P}_r\| = 0\).
Solution

Fix \(r \in [0,1)\). Recall from 11.11 that

\begin{equation*} (\mathcal{P}_r f)(z) = \sum_{n=-\infty}^{\infty} r^{|n|} \hat{f}(n) z^n , \end{equation*}

where the series converges uniformly on \(\partial D\) (with the convention \(0^0 = 1\), so that \(\mathcal{P}_0 f = \hat{f}(0)\), a constant function). Since \(\sigma\) is a probability measure, uniform convergence on \(\partial D\) implies convergence in \(L^2(\partial D)\); hence the displayed series is also the expansion of \(\mathcal{P}_r f\) with respect to the orthonormal basis \(\{z^n\}_{n \in \mathbb{Z}}\) of \(L^2(\partial D)\) (11.30), and therefore

\begin{equation*} \widehat{\mathcal{P}_r f}(n) = r^{|n|} \hat{f}(n) \qquad \text{for all } n \in \mathbb{Z}. \end{equation*}

In other words, \(\mathcal{P}_r\) is the diagonal operator that multiplies the \(n^{\text{th}}\) coordinate by \(r^{|n|}\) in this basis. In particular, by Parseval’s identity [8.63(c)],

\begin{equation*} \|\mathcal{P}_r f\|_2^2 = \sum_{n=-\infty}^{\infty} r^{2|n|} |\hat{f}(n)|^2 \le \sum_{n=-\infty}^{\infty} |\hat{f}(n)|^2 = \|f\|_2^2 , \end{equation*}

so \(\mathcal{P}_r\) is a bounded operator on \(L^2(\partial D)\) with \(\|\mathcal{P}_r\| \le 1\).

(a) Self-adjointness. Suppose \(f, g \in L^2(\partial D)\). Using Parseval’s identity in the inner-product form 8.63(b), with respect to the orthonormal basis \(\{z^n\}_{n \in \mathbb{Z}}\) (11.30), we get

\begin{equation*} \langle \mathcal{P}_r f, g \rangle = \sum_{n=-\infty}^{\infty} \widehat{\mathcal{P}_r f}(n) \overline{\hat{g}(n)} = \sum_{n=-\infty}^{\infty} r^{|n|} \hat{f}(n) \overline{\hat{g}(n)} . \end{equation*}

The same computation applied to \(\langle f, \mathcal{P}_r g \rangle\) gives \(\sum_n \hat{f}(n) \overline{r^{|n|} \hat{g}(n)}\), which equals the sum above because each \(r^{|n|}\) is real. Hence \(\langle \mathcal{P}_r f, g \rangle = \langle f, \mathcal{P}_r g \rangle\) for all \(f, g \in L^2(\partial D)\), so \(\mathcal{P}_r^{*} = \mathcal{P}_r\).

Compactness. For \(M \in \mathbb{Z}^{+}\) define \(T_M\) on \(L^2(\partial D)\) by

\begin{equation*} T_M f = \sum_{|n| \le M} r^{|n|} \hat{f}(n) z^n . \end{equation*}

Then \(T_M\) is bounded and its range lies in the span of \(\{z^n : |n| \le M\}\), which is finite-dimensional; hence \(T_M\) is compact by 10.67. Moreover, for \(f \in L^2(\partial D)\), Parseval’s identity gives

\begin{equation*} \|(\mathcal{P}_r - T_M) f\|_2^2 = \sum_{|n| > M} r^{2|n|} |\hat{f}(n)|^2 \le r^{2(M+1)} \|f\|_2^2 , \end{equation*}

so \(\|\mathcal{P}_r - T_M\| \le r^{M+1}\). Because \(0 \le r < 1\), we have \(r^{M+1} \to 0\) as \(M \to \infty\), so \(\mathcal{P}_r\) is the limit in the norm of \(B(L^2(\partial D))\) of a sequence of compact operators. Since the set of compact operators is closed [10.69(a)], \(\mathcal{P}_r\) is compact.

(b) Suppose \(\lambda \in \mathbb{C}\) and \(f \in L^2(\partial D)\) satisfy \(\mathcal{P}_r f = \lambda f\). Comparing Fourier coefficients gives

\begin{equation*} r^{|n|} \hat{f}(n) = \lambda \hat{f}(n) \qquad \text{for all } n \in \mathbb{Z}, \end{equation*}

so \(\hat{f}(n) = 0\) for every \(n\) with \(r^{|n|} \neq \lambda\). Thus if \(\lambda \notin \{r^{|n|} : n \in \mathbb{Z}\}\), then all Fourier coefficients of \(f\) vanish and hence \(f = 0\) (by 11.43, or by 11.31); such a \(\lambda\) is not an eigenvalue. Conversely each \(r^{|n|}\) is an eigenvalue, with eigenvector \(z^n\).

Case \(r \in (0,1)\). The numbers \(1, r, r^2, r^3, \dots\) are distinct, and \(\{r^{|n|} : n \in \mathbb{Z}\} = \{r^m : m \ge 0\}\). Hence the eigenvalues of \(\mathcal{P}_r\) are exactly \(1, r, r^2, \dots\), and (since \(0\) is not among them and \(\mathcal{P}_r\) has trivial kernel) \(0\) is not an eigenvalue. The eigenspace for the eigenvalue \(1\) consists of those \(f\) with \(\hat{f}(n) = 0\) for all \(n \neq 0\), that is, the constant functions; so the eigenvectors for \(1\) are the nonzero constant functions. For \(m \in \mathbb{Z}^{+}\), the eigenspace for the eigenvalue \(r^m\) consists of those \(f\) with \(\hat{f}(n) = 0\) for all \(n \notin \{m, -m\}\), that is, \(\operatorname{span}\{z^m, z^{-m}\}\); so the eigenvectors for \(r^m\) are the functions \(a z^m + b z^{-m}\) with \(a, b \in \mathbb{C}\) not both \(0\). In particular every nonzero eigenvalue has finite multiplicity, in agreement with 10.82.

Case \(r = 0\). Here \(\mathcal{P}_0 f\) is the constant function whose value is \(\hat{f}(0) = \langle f, 1 \rangle\); that is, \(\mathcal{P}_0 f = \langle f, 1 \rangle 1\), which is the orthogonal projection of \(L^2(\partial D)\) onto the one-dimensional subspace of constant functions. Its eigenvalues are \(1\), with eigenvectors the nonzero constant functions, and \(0\), with eigenvectors the nonzero elements of \(\{f \in L^2(\partial D) : \hat{f}(0) = 0\}\).

(c) The statement is false: in fact \(\|I - \mathcal{P}_r\| = 1\) for every \(r \in [0,1)\), so the limit equals \(1\), not \(0\).

Indeed, \(I - \mathcal{P}_r\) is the diagonal operator with \(\widehat{(I - \mathcal{P}_r)f}(n) = (1 - r^{|n|}) \hat{f}(n)\). Since \(0 \le 1 - r^{|n|} \le 1\) for all \(n\), Parseval’s identity gives \(\|(I - \mathcal{P}_r) f\|_2 \le \|f\|_2\), so \(\|I - \mathcal{P}_r\| \le 1\). On the other hand \(\|z^n\|_2 = 1\) and \(\|(I - \mathcal{P}_r) z^n\|_2 = 1 - r^{|n|}\), which tends to \(1\) as \(n \to \infty\) because \(0 \le r < 1\). Hence \(\|I - \mathcal{P}_r\| \ge \sup_{n} (1 - r^{|n|}) = 1\), and therefore \(\|I - \mathcal{P}_r\| = 1\).

Alternatively, one can argue without computing the norm: if \(\|I - \mathcal{P}_r\| \to 0\) as \(r \uparrow 1\), then \(I\) would be a limit in operator norm of the compact operators \(\mathcal{P}_r\) and hence compact by 10.69(a). But the identity operator on the infinite-dimensional Hilbert space \(L^2(\partial D)\) is not compact, since the orthonormal sequence \(z, z^2, z^3, \dots\) is bounded and has no convergent subsequence (any two distinct terms are at distance \(\sqrt{2}\)).

Problem (11B.10)

Suppose \(f \in L^1(\partial D)\). Define \(T \colon L^2(\partial D) \to L^2(\partial D)\) by \(Tg = f * g\).

  • (a) Show that \(T\) is a compact operator on \(L^2(\partial D)\).
  • (b) Prove that \(T\) is injective if and only if \(\hat{f}(n) \neq 0\) for every \(n \in \mathbb{Z}\).
  • (c) Find a formula for \(T^{*}\).
  • (d) Prove: \(T\) is self-adjoint if and only if all Fourier coefficients of \(f\) are real.
  • (e) Show that \(T\) is a normal operator.
Solution

First note that \(T\) really does map \(L^2(\partial D)\) into \(L^2(\partial D)\): because \(\sigma\) is a finite measure we have \(L^2(\partial D) \subseteq L^1(\partial D)\), so \(f * g\) is defined almost everywhere and lies in \(L^1(\partial D)\) by 11.37, and 11.38 (with \(p = 2\)) gives

\begin{equation*} \|Tg\|_2 = \|f * g\|_2 \le \|f\|_1 \|g\|_2 . \end{equation*}

Thus \(T\) is a bounded operator on \(L^2(\partial D)\) with \(\|T\| \le \|f\|_1\); linearity of \(T\) is immediate from the definition of convolution. The key structural fact is 11.44:

\begin{equation*} \widehat{Tg}(n) = \widehat{f * g}(n) = \hat{f}(n) \hat{g}(n) \qquad \text{for all } n \in \mathbb{Z}, \; g \in L^2(\partial D). \end{equation*}

So in the orthonormal basis \(\{z^n\}_{n \in \mathbb{Z}}\) of \(L^2(\partial D)\) (11.30), \(T\) is the diagonal operator whose \(n^{\text{th}}\) multiplier is \(\hat{f}(n)\). Throughout we use Parseval’s identity [8.63(c)]: \(\|h\|_2^2 = \sum_{n} |\hat{h}(n)|^2\) for \(h \in L^2(\partial D)\), and the fact that an element of \(L^2(\partial D)\) is determined by its Fourier coefficients (11.43).

(a) By the Riemann–Lebesgue Lemma (11.10), \(\lim_{n \to \pm \infty} \hat{f}(n) = 0\). For \(M \in \mathbb{Z}^{+}\) let \(Q_M\) denote the orthogonal projection of \(L^2(\partial D)\) onto \(\operatorname{span}\{z^n : |n| \le M\}\), and set \(T_M = Q_M T\), so that

\begin{equation*} T_M g = \sum_{|n| \le M} \hat{f}(n) \hat{g}(n) z^n . \end{equation*}

Each \(T_M\) is a bounded operator whose range is contained in the finite-dimensional space \(\operatorname{span}\{z^n : |n| \le M\}\), so \(T_M\) is compact by 10.67. Set \(\varepsilon_M = \sup_{|n| > M} |\hat{f}(n)|\); by the Riemann–Lebesgue Lemma, \(\varepsilon_M \to 0\) as \(M \to \infty\). For \(g \in L^2(\partial D)\), Parseval’s identity gives

\begin{equation*} \|(T - T_M) g\|_2^2 = \sum_{|n| > M} |\hat{f}(n)|^2 |\hat{g}(n)|^2 \le \varepsilon_M^2 \|g\|_2^2 , \end{equation*}

hence \(\|T - T_M\| \le \varepsilon_M \to 0\). Because the set of compact operators on a Hilbert space is closed in the operator norm [10.69(a)], \(T\) is compact.

(b) Suppose \(\hat{f}(m) = 0\) for some \(m \in \mathbb{Z}\). Then \(\widehat{T z^m}(n) = \hat{f}(n) \widehat{z^m}(n)\), which is \(0\) for \(n \neq m\) (because \(\widehat{z^m}(n) = 0\) there, by 11.6) and is \(\hat{f}(m) = 0\) for \(n = m\). So all Fourier coefficients of \(T z^m\) vanish, whence \(T z^m = 0\) by 11.43. Since \(z^m \neq 0\), the operator \(T\) is not injective.

Conversely, suppose \(\hat{f}(n) \neq 0\) for every \(n \in \mathbb{Z}\), and suppose \(g \in L^2(\partial D)\) with \(Tg = 0\). Then \(\hat{f}(n) \hat{g}(n) = 0\) for every \(n\), so \(\hat{g}(n) = 0\) for every \(n\), so \(g = 0\) by 11.43. Hence \(T\) is injective.

(c) Define \(f^{\#} \colon \partial D \to \mathbb{C}\) by \(f^{\#}(z) = \overline{f(\bar{z})}\). The map \(z \mapsto \bar{z}\) is a homeomorphism of \(\partial D\) that corresponds under 11.2 to \(t \mapsto -t\) on \((-\pi, \pi]\), so it carries measurable subsets of \(\partial D\) to measurable subsets; hence \(f^{\#}\) is measurable. Taking \(z = 1\) in 11.17 gives

\begin{equation*} \int_{\partial D} h(\bar{w})\, d\sigma(w) = \int_{\partial D} h(\zeta)\, d\sigma(\zeta) \end{equation*}

for every measurable \(h \colon \partial D \to \mathbb{C}\) for which either side is defined. Applying this to \(h = |f|\) shows \(f^{\#} \in L^1(\partial D)\) with \(\|f^{\#}\|_1 = \|f\|_1\), and applying it to \(h(\zeta) = \overline{f(\zeta)}\, \zeta^n\) gives

\begin{equation*} \widehat{f^{\#}}(n) = \int_{\partial D} \overline{f(\bar{z})}\, \overline{z}^{\,n}\, d\sigma(z) = \int_{\partial D} \overline{f(w)}\, w^n \, d\sigma(w) = \overline{\int_{\partial D} f(w)\, \overline{w}^{\,n}\, d\sigma(w)} = \overline{\hat{f}(n)} , \end{equation*}

where the first equality used \(\overline{z}^{\,n} = w^n\) for \(w = \bar{z}\).

I claim that \(T^{*} h = f^{\#} * h\). Indeed, for \(g, h \in L^2(\partial D)\), using Parseval’s identity in the inner-product form 8.63(b) with respect to the orthonormal basis \(\{z^n\}_{n \in \mathbb{Z}}\) (11.30), together with 11.44 twice,

\begin{equation*} \langle Tg, h \rangle = \sum_{n=-\infty}^{\infty} \hat{f}(n) \hat{g}(n) \overline{\hat{h}(n)} = \sum_{n=-\infty}^{\infty} \hat{g}(n) \overline{\overline{\hat{f}(n)} \hat{h}(n)} = \sum_{n=-\infty}^{\infty} \hat{g}(n) \overline{\widehat{f^{\#} * h}(n)} = \langle g, f^{\#} * h \rangle ; \end{equation*}

all these sums converge absolutely because \(\{\hat{f}(n)\}\) is bounded [11.9(c)] while \(\{\hat{g}(n)\}\) and \(\{\hat{h}(n)\}\) are in \(\ell^2(\mathbb{Z})\). Since \(h \mapsto f^{\#} * h\) is a bounded operator on \(L^2(\partial D)\) (by the first paragraph, applied to \(f^{\#}\)), the uniqueness of the adjoint gives

\begin{equation*} T^{*} h = f^{\#} * h, \qquad \text{where } f^{\#}(z) = \overline{f(\bar{z})} . \end{equation*}

Written out as an integral and simplified by another application of 11.17 with \(z = 1\), this says

\begin{equation*} (T^{*} h)(z) = \int_{\partial D} \overline{f(\bar{w})}\, h(z \bar{w})\, d\sigma(w) = \int_{\partial D} \overline{f(w)}\, h(zw)\, d\sigma(w) . \end{equation*}

Equivalently, \(\widehat{T^{*} h}(n) = \overline{\hat{f}(n)} \hat{h}(n)\) for all \(n \in \mathbb{Z}\): \(T^{*}\) is the diagonal operator with multipliers \(\overline{\hat{f}(n)}\).

(d) Suppose \(T\) is self-adjoint. Fix \(m \in \mathbb{Z}\) and apply \(T = T^{*}\) to \(z^m\). Since \(\widehat{z^m}(m) = 1\) by 11.6, comparing \(m^{\text{th}}\) Fourier coefficients and using the last display of (c) gives \(\hat{f}(m) = \overline{\hat{f}(m)}\), so \(\hat{f}(m) \in \mathbb{R}\). As \(m\) was arbitrary, all Fourier coefficients of \(f\) are real.

Conversely, suppose \(\hat{f}(n) \in \mathbb{R}\) for all \(n \in \mathbb{Z}\). Then \(\widehat{f^{\#}}(n) = \overline{\hat{f}(n)} = \hat{f}(n)\) for all \(n\), so \(f^{\#} = f\) by 11.43, and hence \(T^{*} h = f^{\#} * h = f * h = Th\) for all \(h \in L^2(\partial D)\). Thus \(T\) is self-adjoint.

(e) For every \(g \in L^2(\partial D)\) and every \(n \in \mathbb{Z}\), part (c) and 11.44 give

\begin{equation*} \widehat{TT^{*} g}(n) = \hat{f}(n) \overline{\hat{f}(n)} \hat{g}(n) = |\hat{f}(n)|^2 \hat{g}(n) = \overline{\hat{f}(n)} \hat{f}(n) \hat{g}(n) = \widehat{T^{*} T g}(n) . \end{equation*}

Thus \(TT^{*} g\) and \(T^{*} T g\) have the same Fourier coefficients, so \(TT^{*} g = T^{*} T g\) by 11.43. Since \(g\) was arbitrary, \(TT^{*} = T^{*} T\), which means \(T\) is normal.

Problem (11B.11)

Show that if \(f, g \in L^1(\partial D)\) then

\begin{equation*} (f * g)^{\sim}(t) = \frac{1}{2\pi} \int_{-\pi}^{\pi} \tilde{f}(x)\, \tilde{g}(t - x)\, dx, \end{equation*}

for those \(t \in \mathbb{R}\) such that \((f * g)(e^{it})\) makes sense; here \((f * g)^{\sim}\), \(\tilde{f}\), and \(\tilde{g}\) denote the transfers to the real line as defined in 11.24.

Solution

Recall from 11.24 that for a function \(h \colon \partial D \to \mathbb{C}\) the transfer \(\tilde{h} \colon \mathbb{R} \to \mathbb{C}\) is defined by \(\tilde{h}(t) = h(e^{it})\); note that \(\tilde{h}\) is \(2\pi\)-periodic, so \(\tilde{h}(s)\) makes sense for every real \(s\) regardless of whether \(s\) lies in \((-\pi, \pi]\).

Recall also the transfer of integration from \(\partial D\) to \((-\pi, \pi]\) recorded after 11.4: for every measurable \(h \colon \partial D \to \mathbb{C}\) for which either side is defined,

\begin{equation*} \int_{\partial D} h \, d\sigma = \int_{-\pi}^{\pi} h(e^{ix}) \, \frac{dx}{2\pi} . \end{equation*}

Now fix \(t \in \mathbb{R}\) and put \(z = e^{it}\). By the definition of convolution (11.36),

\begin{equation*} (f * g)(z) = \int_{\partial D} f(w)\, g(z \bar{w}) \, d\sigma(w), \end{equation*}

and by hypothesis \(t\) is such that this integral makes sense. The function \(h \colon w \mapsto f(w) g(z \bar{w})\) is measurable: \(w \mapsto z \bar{w}\) corresponds under 11.2 to \(x \mapsto t - x\) on \((-\pi, \pi]\) (modulo \(2\pi\)), which is a Borel isomorphism, so \(w \mapsto g(z \bar{w})\) is measurable, and a product of measurable functions is measurable.

Apply the transfer of integration to \(h(w) = f(w) g(z \bar{w})\). Writing \(w = e^{ix}\) with \(x \in (-\pi, \pi]\), we have \(\bar{w} = e^{-ix}\) and hence

\begin{equation*} z \bar{w} = e^{it} e^{-ix} = e^{i(t - x)} . \end{equation*}

Therefore \(h(e^{ix}) = f(e^{ix}) g(e^{i(t-x)}) = \tilde{f}(x)\, \tilde{g}(t - x)\), and the transfer of integration gives

\begin{equation*} (f * g)(e^{it}) = \int_{\partial D} f(w)\, g(z\bar{w})\, d\sigma(w) = \int_{-\pi}^{\pi} \tilde{f}(x)\, \tilde{g}(t-x)\, \frac{dx}{2\pi} = \frac{1}{2\pi} \int_{-\pi}^{\pi} \tilde{f}(x)\, \tilde{g}(t-x)\, dx . \end{equation*}

Since \((f * g)^{\sim}(t)\) is by definition \((f * g)(e^{it})\), this is exactly the asserted identity; in particular the integral on the right is defined for precisely those \(t\) for which \((f * g)(e^{it})\) makes sense, which by 11.37 is almost every \(t \in \mathbb{R}\).

Thus convolution on \(\partial D\) corresponds, under the transfer to \(\mathbb{R}\), to the usual convolution of \(2\pi\)-periodic functions over a period, normalized by \(\frac{1}{2\pi}\).

Problem (11B.12)

Suppose \(1 \le p \le \infty\). Prove that if \(f \in L^p(\partial D)\) and \(g \in L^{p^{\prime}}(\partial D)\), then \(f * g\) is a continuous function on \(\partial D\).

Solution

Here \(p^{\prime}\) denotes the exponent conjugate to \(p\), so that \(\frac{1}{p} + \frac{1}{p^{\prime}} = 1\) (with \(1^{\prime} = \infty\) and \(\infty^{\prime} = 1\)). For \(h \colon \partial D \to \mathbb{C}\) and \(z \in \partial D\), write \(h_z\) for the rotated-and-reflected function

\begin{equation*} h_z(w) = h(z \bar{w}) \qquad (w \in \partial D). \end{equation*}

Because \(w \mapsto z \bar{w}\) is a homeomorphism of \(\partial D\) that preserves \(\sigma\) (this is the invariance recorded in 11.17), \(h_z\) is measurable whenever \(h\) is, and for every \(q \in [1, \infty]\),

\begin{equation*} \|h_z\|_q = \|h\|_q . \end{equation*}

Indeed for \(q < \infty\) this follows by applying 11.17 to \(|h|^q\), and for \(q = \infty\) it holds because \(z \bar{\cdot}\) maps sets of measure zero to sets of measure zero.

Step 1: \((f * g)(z)\) is defined at every point of \(\partial D\).

Let \(z \in \partial D\). By Hölder’s inequality (7.9) and the displayed norm identity,

\begin{equation*} \int_{\partial D} |f(w) g(z \bar{w})| \, d\sigma(w) \le \|f\|_p \, \|g_z\|_{p^{\prime}} = \|f\|_p \, \|g\|_{p^{\prime}} < \infty . \end{equation*}

Hence the integral defining \((f * g)(z)\) converges absolutely for every \(z \in \partial D\), and \(|(f*g)(z)| \le \|f\|_p \|g\|_{p^{\prime}}\).

Step 2: continuity of translation in \(L^q\) for \(q < \infty\).

Claim: if \(1 \le q < \infty\) and \(h \in L^q(\partial D)\), then the map \(z \mapsto h_z\) is continuous from \(\partial D\) to \(L^q(\partial D)\).

To prove this, let \(z_0 \in \partial D\) and let \(\varepsilon > 0\). Because \(q < \infty\), the continuous functions are dense in \(L^q(\partial D)\) [this is the approximation used in the proof of 11.42; it follows from 3.47 and 3.48 with the \(L^1\)-norm replaced by the \(L^q\)-norm, exactly as explained in the bracketed remark in the proof of 11.30]. So choose a continuous \(\varphi \colon \partial D \to \mathbb{C}\) with \(\|h - \varphi\|_q < \varepsilon\). Then for every \(z \in \partial D\),

\begin{equation*} \|h_z - h_{z_0}\|_q \le \|h_z - \varphi_z\|_q + \|\varphi_z - \varphi_{z_0}\|_q + \|\varphi_{z_0} - h_{z_0}\|_q < 2\varepsilon + \|\varphi_z - \varphi_{z_0}\|_q , \end{equation*}

where we used \(\|h_z - \varphi_z\|_q = \|(h - \varphi)_z\|_q = \|h - \varphi\|_q\) and similarly for \(z_0\). Now \(\partial D\) is compact, so \(\varphi\) is uniformly continuous; and for every \(w \in \partial D\),

\begin{equation*} |z \bar{w} - z_0 \bar{w}| = |z - z_0| \, |\bar{w}| = |z - z_0| . \end{equation*}

Hence \(\sup_{w \in \partial D} |\varphi_z(w) - \varphi_{z_0}(w)| \to 0\) as \(z \to z_0\), and since \(\sigma(\partial D) = 1\) we get \(\|\varphi_z - \varphi_{z_0}\|_q \le \|\varphi_z - \varphi_{z_0}\|_\infty \to 0\). Therefore \(\limsup_{z \to z_0} \|h_z - h_{z_0}\|_q \le 2\varepsilon\), and as \(\varepsilon > 0\) was arbitrary, \(\lim_{z \to z_0} \|h_z - h_{z_0}\|_q = 0\). This proves the claim.

Step 3: \(f * g\) is continuous.

At least one of \(p, p^{\prime}\) is finite. Suppose first \(p^{\prime} < \infty\). Fix \(z_0 \in \partial D\). For \(z \in \partial D\), Hölder’s inequality (7.9) gives

\begin{equation*} |(f * g)(z) - (f * g)(z_0)| = \left| \int_{\partial D} f(w) \big( g_z(w) - g_{z_0}(w) \big) d\sigma(w) \right| \le \|f\|_p \, \|g_z - g_{z_0}\|_{p^{\prime}} , \end{equation*}

and the right side tends to \(0\) as \(z \to z_0\) by Step 2 applied to \(h = g \in L^{p^{\prime}}(\partial D)\) with \(q = p^{\prime}\). Hence \(f * g\) is continuous at \(z_0\), and since \(z_0\) was arbitrary, \(f * g\) is continuous on \(\partial D\).

If instead \(p^{\prime} = \infty\), then \(p = 1 < \infty\), and we interchange the roles of the two functions. By 11.41 we have \(f * g = g * f\), so for \(z, z_0 \in \partial D\),

\begin{equation*} |(g * f)(z) - (g * f)(z_0)| = \left| \int_{\partial D} g(w) \big( f_z(w) - f_{z_0}(w) \big) d\sigma(w) \right| \le \|g\|_{\infty} \, \|f_z - f_{z_0}\|_1 , \end{equation*}

which tends to \(0\) as \(z \to z_0\) by Step 2 applied to \(h = f \in L^1(\partial D)\) with \(q = 1\). So \(g * f\) is continuous, and \(f * g = g * f\) is continuous on \(\partial D\).

In both cases \(f * g\) is continuous on \(\partial D\), as desired. (The commutativity 11.41 was stated for functions in \(L^1(\partial D)\), which covers the present situation: since \(\sigma\) is a finite measure, \(L^q(\partial D) \subseteq L^1(\partial D)\) for every \(q \in [1, \infty]\).)

Problem (11B.13)

Suppose \(g \in L^1(\partial D)\) is such that \(\hat{g}(n) \neq 0\) for infinitely many \(n \in \mathbb{Z}\). Prove that if \(f \in L^1(\partial D)\), then \(f * g \neq g\).

Solution

Suppose, for contradiction, that there exists \(f \in L^1(\partial D)\) with \(f * g = g\) (as elements of \(L^1(\partial D)\), i.e. equality almost everywhere; recall from 11.37 that \(f * g\) is defined almost everywhere and lies in \(L^1(\partial D)\)).

Taking Fourier coefficients and using 11.44, we get

\begin{equation*} \hat{f}(n) \hat{g}(n) = \widehat{f * g}(n) = \hat{g}(n) \qquad \text{for every } n \in \mathbb{Z}, \end{equation*}

that is, \(\big( \hat{f}(n) - 1 \big) \hat{g}(n) = 0\) for every \(n \in \mathbb{Z}\).

Let \(S = \{n \in \mathbb{Z} : \hat{g}(n) \neq 0\}\). By hypothesis \(S\) is infinite, and the displayed equation forces

\begin{equation*} \hat{f}(n) = 1 \qquad \text{for every } n \in S . \end{equation*}

Since \(S\) is an infinite subset of \(\mathbb{Z}\), it is unbounded: for each \(M \in \mathbb{Z}^{+}\) the set \(\{n \in \mathbb{Z} : |n| \le M\}\) is finite, so \(S\) contains an element \(n\) with \(|n| > M\). Hence we may choose \(n_1, n_2, \dots \in S\) with \(|n_k| \to \infty\), and after passing to a subsequence we may assume that \(n_k \to \infty\) or that \(n_k \to -\infty\).

But the Riemann–Lebesgue Lemma (11.10) applied to \(f \in L^1(\partial D)\) says that \(\lim_{n \to \pm \infty} \hat{f}(n) = 0\), so \(\hat{f}(n_k) \to 0\). This contradicts \(\hat{f}(n_k) = 1\) for every \(k\).

Therefore no such \(f\) exists; in other words, \(f * g \neq g\) for every \(f \in L^1(\partial D)\).

Remark: such a \(g\) exists. For instance, the second bullet point of Example 11.8 exhibits \(g(z) = \frac{1}{|3 - z|^2}\), which is continuous on \(\partial D\) and satisfies \(\hat{g}(n) = \frac{1}{8} \cdot \frac{1}{3^{|n|}} \neq 0\) for every \(n \in \mathbb{Z}\). Hence the result just proved shows that \(L^1(\partial D)\), with convolution as multiplication, has no multiplicative identity: there is no \(f \in L^1(\partial D)\) with \(f * g = g\) for all \(g \in L^1(\partial D)\). This is the reason one works instead with approximate identities such as the Poisson kernel \(\{P_r\}_{r \in [0,1)}\), whose defining properties are recorded in 11.16.

Problem (11B.14)

Show that there exists a two-sided sequence \(\dots, b_{-2}, b_{-1}, b_0, b_1, b_2, \dots\) such that \(\lim_{n \to \pm \infty} b_n = 0\) but there does not exist \(f \in L^1(\partial D)\) with \(\hat{f}(n) = b_n\) for all \(n \in \mathbb{Z}\).

Solution

In other words, we must show that the converse of the Riemann–Lebesgue Lemma (11.10) fails: not every two-sided sequence tending to \(0\) at \(\pm \infty\) is the sequence of Fourier coefficients of a function in \(L^1(\partial D)\).

Setup. Let

\begin{equation*} c_0(\mathbb{Z}) = \Big\{ b = (b_n)_{n \in \mathbb{Z}} : b_n \in \mathbb{C} \text{ and } \lim_{n \to \pm \infty} b_n = 0 \Big\}, \qquad \|b\|_{\infty} = \sup_{n \in \mathbb{Z}} |b_n| . \end{equation*}

Then \(c_0(\mathbb{Z})\) is a Banach space: it is a subspace of \(\ell^{\infty}(\mathbb{Z})\), which is \(L^{\infty}\) of counting measure on \(\mathbb{Z}\) and hence a Banach space by 7.24, and \(c_0(\mathbb{Z})\) is closed in \(\ell^{\infty}(\mathbb{Z})\) by the usual uniform-limit argument: if \(b^{(k)} \in c_0(\mathbb{Z})\) and \(\|b^{(k)} - b\|_{\infty} \to 0\), then given \(\varepsilon > 0\) pick \(k\) with \(\|b^{(k)} - b\|_{\infty} < \varepsilon/2\) and then \(M\) with \(|b^{(k)}_n| < \varepsilon/2\) for \(|n| \ge M\), so that \(|b_n| < \varepsilon\) for \(|n| \ge M\).

Define \(\Phi \colon L^1(\partial D) \to c_0(\mathbb{Z})\) by \(\Phi(f) = \big( \hat{f}(n) \big)_{n \in \mathbb{Z}}\). By the Riemann–Lebesgue Lemma (11.10), \(\Phi\) does map into \(c_0(\mathbb{Z})\); by 11.9(a) and 11.9(b) it is linear; by 11.9(c) we have \(\|\Phi(f)\|_{\infty} \le \|f\|_1\), so \(\Phi\) is bounded; and by 11.43 it is injective. Also \(L^1(\partial D)\) is a Banach space (7.24).

The exercise asks us to show that \(\Phi\) is not surjective. Suppose, for contradiction, that \(\Phi\) is surjective. Then \(\Phi\) is a one-to-one bounded linear map from the Banach space \(L^1(\partial D)\) onto the Banach space \(c_0(\mathbb{Z})\), so the Bounded Inverse Theorem (6.83) says \(\Phi^{-1}\) is bounded. Thus there exists \(c \in (0, \infty)\) such that

\begin{equation*} \|f\|_1 \le c \, \sup_{n \in \mathbb{Z}} |\hat{f}(n)| \qquad \text{for every } f \in L^1(\partial D). \end{equation*}

The Dirichlet kernel contradicts this. For \(M \in \mathbb{Z}^{+}\) define \(D_M \colon \partial D \to \mathbb{C}\) by

\begin{equation*} D_M(z) = \sum_{n = -M}^{M} z^n . \end{equation*}

Each \(D_M\) is continuous, hence in \(L^1(\partial D)\), and by 11.6 its Fourier coefficients are

\begin{equation*} \widehat{D_M}(n) = \begin{cases} 1 & \text{if } |n| \le M,\\ 0 & \text{if } |n| > M,\end{cases} \end{equation*}

so \(\sup_{n} |\widehat{D_M}(n)| = 1\). Hence the inequality above would give \(\|D_M\|_1 \le c\) for every \(M\). We show this is false by proving \(\|D_M\|_1 \to \infty\).

Summing the geometric series, for \(t \in \mathbb{R}\) with \(e^{it} \neq 1\),

\begin{equation*} D_M(e^{it}) = \sum_{n=-M}^{M} e^{int} = \frac{e^{i(M + \frac{1}{2})t} - e^{-i(M + \frac{1}{2})t}}{e^{it/2} - e^{-it/2}} = \frac{\sin\big( (M + \frac{1}{2}) t \big)}{\sin \frac{t}{2}} . \end{equation*}

Therefore, transferring the integral to \((-\pi, \pi]\) as after 11.4,

\begin{equation*} \|D_M\|_1 = \int_{-\pi}^{\pi} \left| \frac{\sin\big( (M + \frac{1}{2}) t \big)}{\sin \frac{t}{2}} \right| \frac{dt}{2\pi} . \end{equation*}

Since \(|\sin s| \le |s|\) for all real \(s\), we have \(\big| \sin \frac{t}{2} \big| \le \frac{|t|}{2}\), and so

\begin{equation*} \|D_M\|_1 \ge \frac{1}{2\pi} \int_{-\pi}^{\pi} \frac{2 \big| \sin\big( (M + \frac{1}{2}) t \big) \big|}{|t|} \, dt = \frac{2}{\pi} \int_{0}^{\pi} \frac{\big| \sin\big( (M + \frac{1}{2}) t \big) \big|}{t} \, dt , \end{equation*}

using that the integrand is even. Substituting \(s = (M + \frac{1}{2}) t\) (so \(\frac{dt}{t} = \frac{ds}{s}\)) gives

\begin{equation*} \|D_M\|_1 \ge \frac{2}{\pi} \int_{0}^{(M + \frac{1}{2}) \pi} \frac{|\sin s|}{s} \, ds \ge \frac{2}{\pi} \int_{0}^{M\pi} \frac{|\sin s|}{s} \, ds = \frac{2}{\pi} \sum_{k=1}^{M} \int_{(k-1)\pi}^{k\pi} \frac{|\sin s|}{s} \, ds . \end{equation*}

On \([(k-1)\pi, k\pi]\) we have \(s \le k\pi\) and \(\int_{(k-1)\pi}^{k\pi} |\sin s| \, ds = 2\), so

\begin{equation*} \|D_M\|_1 \ge \frac{2}{\pi} \sum_{k=1}^{M} \frac{2}{k \pi} = \frac{4}{\pi^2} \sum_{k=1}^{M} \frac{1}{k} \longrightarrow \infty \quad \text{as } M \to \infty , \end{equation*}

because the harmonic series diverges. (This divergence is also the content of Exercise 11(b) in Section 11A.)

This contradicts \(\|D_M\|_1 \le c\) for all \(M\). Hence \(\Phi\) is not surjective: there exists \(b = (b_n)_{n \in \mathbb{Z}} \in c_0(\mathbb{Z})\), that is, a two-sided sequence with \(\lim_{n \to \pm \infty} b_n = 0\), such that no \(f \in L^1(\partial D)\) satisfies \(\hat{f}(n) = b_n\) for all \(n \in \mathbb{Z}\).

Problem (11B.15)

Prove that if \(f, g \in L^2(\partial D)\), then

\begin{equation*} \widehat{fg}(n) = \sum_{k=-\infty}^{\infty} \hat{f}(k)\, \hat{g}(n-k) \end{equation*}

for every \(n \in \mathbb{Z}\).

Solution

Throughout, \(\langle u, v \rangle = \int_{\partial D} u \bar{v} \, d\sigma\) denotes the inner product of \(L^2(\partial D)\), and \(\hat{u}(n) = \int_{\partial D} u(z)\, \bar{z}^{\,n} \, d\sigma(z)\) denotes the \(n^{\text{th}}\) Fourier coefficient of \(u \in L^1(\partial D)\); when \(u \in L^2(\partial D)\) this reads \(\hat{u}(n) = \langle u, z^n \rangle\). Note \(L^2(\partial D) \subseteq L^1(\partial D)\) because \(\sigma\) is a finite measure, so all Fourier coefficients below are defined.

First note that both sides of the asserted identity make sense. By Hölder’s inequality (7.9), \(\|fg\|_1 \le \|f\|_2 \|g\|_2 < \infty\), so \(fg \in L^1(\partial D)\) and \(\widehat{fg}(n)\) is defined. Also, \(\{z^k\}_{k \in \mathbb{Z}}\) is an orthonormal family in \(L^2(\partial D)\) (11.6), so Bessel’s inequality (8.57) gives \(\sum_{k \in \mathbb{Z}} |\hat{f}(k)|^2 \le \|f\|_2^2\) and \(\sum_{k \in \mathbb{Z}} |\hat{g}(k)|^2 \le \|g\|_2^2\). Hence, by the Cauchy-Schwarz inequality in \(\ell^2(\mathbb{Z})\),

\begin{equation*} \sum_{k=-\infty}^{\infty} \bigl|\hat{f}(k)\bigr| \, \bigl|\hat{g}(n-k)\bigr| \le \Bigl( \sum_{k \in \mathbb{Z}} |\hat{f}(k)|^2 \Bigr)^{1/2} \Bigl( \sum_{k \in \mathbb{Z}} |\hat{g}(n-k)|^2 \Bigr)^{1/2} \le \|f\|_2 \|g\|_2 < \infty , \end{equation*}

where we used that \(k \mapsto n-k\) is a bijection of \(\mathbb{Z}\) onto itself. Thus the series on the right side converges absolutely, and in particular its value does not depend on the order of summation.

Now fix \(n \in \mathbb{Z}\) and define \(h \colon \partial D \to \mathbb{C}\) by

\begin{equation*} h(z) = \overline{g(z)}\, z^n . \end{equation*}

Because \(|h| = |g|\) pointwise, we have \(h \in L^2(\partial D)\) with \(\|h\|_2 = \|g\|_2\). The point of this definition is that for every \(z \in \partial D\) we have \(\overline{h(z)} = g(z) \bar{z}^{\,n}\) (using \(\bar z^{\,n} = \overline{z^n}\)), and therefore

\begin{equation*} \widehat{fg}(n) = \int_{\partial D} f(z) g(z) \bar{z}^{\,n} \, d\sigma(z) = \int_{\partial D} f(z) \overline{h(z)} \, d\sigma(z) = \langle f, h \rangle . \end{equation*}

Next we compute the Fourier coefficients of \(h\). For \(k \in \mathbb{Z}\),

\begin{equation*} \hat{h}(k) = \int_{\partial D} \overline{g(z)}\, z^n \bar{z}^{\,k} \, d\sigma(z) = \int_{\partial D} \overline{\, g(z)\, \bar{z}^{\,n} z^{k} \,} \, d\sigma(z) = \overline{\int_{\partial D} g(z)\, \bar{z}^{\,n-k} \, d\sigma(z)} = \overline{\hat{g}(n-k)} , \end{equation*}

where the middle equality is the algebraic identity \(\overline{\,g(z)\, \bar z^{\,n} z^{k}\,} = \overline{g(z)}\, z^n \bar z^{\,k}\), the third equality pulls the conjugation out of the integral and uses \(\bar z^{\,n} z^{k} = \bar z^{\,n-k}\) (valid because \(z \bar z = |z|^2 = 1\) for \(z \in \partial D\)), and the last equality is the definition of \(\hat{g}(n-k)\). Equivalently, \(\overline{\hat{h}(k)} = \hat{g}(n-k)\).

Finally we expand \(\langle f, h \rangle\). Since \(\{z^k\}_{k \in \mathbb{Z}}\) is an orthonormal basis of \(L^2(\partial D)\) (11.30), Parseval’s identity [8.63(b)] gives

\begin{equation*} \langle f, h \rangle = \sum_{k=-\infty}^{\infty} \langle f, z^k \rangle \, \overline{\langle h, z^k \rangle} = \sum_{k=-\infty}^{\infty} \hat{f}(k)\, \overline{\hat{h}(k)} = \sum_{k=-\infty}^{\infty} \hat{f}(k)\, \hat{g}(n-k) . \end{equation*}

Combining the last display with \(\widehat{fg}(n) = \langle f, h \rangle\) completes the proof.

If one prefers to avoid quoting 8.63(b), the same conclusion follows directly from 11.31: the partial sums \(S_K = \sum_{|k| \le K} \hat{f}(k) z^k\) converge to \(f\) in the norm of \(L^2(\partial D)\), and the inner product is continuous in its first slot (by the Cauchy-Schwarz inequality 8.11), so

\begin{equation*} \langle f, h \rangle = \lim_{K \to \infty} \langle S_K, h \rangle = \lim_{K \to \infty} \sum_{|k| \le K} \hat{f}(k) \overline{\hat{h}(k)} = \sum_{k=-\infty}^{\infty} \hat{f}(k)\, \hat{g}(n-k), \end{equation*}

the last equality holding because the series converges absolutely, as shown above.

Problem (11B.16)

Suppose \(f \in L^1(\partial D)\). Prove that \(\mathcal{P}_r(\mathcal{P}_s f) = \mathcal{P}_{rs} f\) for all \(r, s \in [0,1)\).

Solution

Recall from 11.11 that for \(u \in L^1(\partial D)\) and \(t \in [0,1)\),

\begin{equation*} (\mathcal{P}_t u)(z) = \sum_{n=-\infty}^{\infty} t^{|n|} \hat{u}(n) z^n \qquad (z \in \partial D), \end{equation*}

with the convention \(0^0 = 1\), and that the series converges uniformly on \(\partial D\); in particular \(\mathcal{P}_t u\) is continuous on \(\partial D\), hence lies in \(L^1(\partial D)\), so that \(\mathcal{P}_r(\mathcal{P}_s f)\) is defined.

Step 1: the Fourier coefficients of \(\mathcal{P}_s f\). Fix \(s \in [0,1)\) and \(n \in \mathbb{Z}\). As noted in 11.11, the estimate \(|s^{|m|} \hat{f}(m) z^m| \le \|f\|_1 s^{|m|}\) for \(z \in \partial D\) together with \(\sum_{m \in \mathbb{Z}} s^{|m|} < \infty\) shows that the series \(\sum_{m} s^{|m|} \hat{f}(m) z^m\) converges uniformly on \(\partial D\). Since \(\sigma(\partial D) = 1\) and \(|\bar{z}^{\,n}| = 1\) on \(\partial D\), uniform convergence permits termwise integration against \(\bar{z}^{\,n}\):

\begin{equation*} \widehat{\mathcal{P}_s f}(n) = \int_{\partial D} (\mathcal{P}_s f)(z)\, \bar{z}^{\,n} \, d\sigma(z) = \sum_{m=-\infty}^{\infty} s^{|m|} \hat{f}(m) \int_{\partial D} z^m \bar{z}^{\,n} \, d\sigma(z) . \end{equation*}

Because \(\{z^m\}_{m \in \mathbb{Z}}\) is an orthonormal family in \(L^2(\partial D)\) (11.6), we have \(\int_{\partial D} z^m \bar{z}^{\,n} \, d\sigma = 1\) if \(m = n\) and \(0\) otherwise. Hence

\begin{equation*} \widehat{\mathcal{P}_s f}(n) = s^{|n|} \hat{f}(n) \qquad \text{for every } n \in \mathbb{Z}. \end{equation*}

Step 2: conclusion. Now let \(r, s \in [0,1)\) and let \(z \in \partial D\). Applying the definition 11.11 to the function \(\mathcal{P}_s f \in L^1(\partial D)\) and then using Step 1,

\begin{equation*} \bigl(\mathcal{P}_r(\mathcal{P}_s f)\bigr)(z) = \sum_{n=-\infty}^{\infty} r^{|n|} \widehat{\mathcal{P}_s f}(n) z^n = \sum_{n=-\infty}^{\infty} r^{|n|} s^{|n|} \hat{f}(n) z^n = \sum_{n=-\infty}^{\infty} (rs)^{|n|} \hat{f}(n) z^n = (\mathcal{P}_{rs} f)(z), \end{equation*}

where the third equality uses \(r^{|n|} s^{|n|} = (rs)^{|n|}\) and the last equality is the definition 11.11 applied with \(rs \in [0,1)\) in place of \(r\). All series here converge absolutely, so the manipulations are legitimate.

Since \(z \in \partial D\) was arbitrary, \(\mathcal{P}_r(\mathcal{P}_s f) = \mathcal{P}_{rs} f\), as desired.

Note that the convention \(0^0 = 1\) is consistent with this computation: if, say, \(r = 0\), then both sides reduce to the constant function \(\hat{f}(0)\), because \(\widehat{\mathcal{P}_s f}(0) = s^0 \hat{f}(0) = \hat{f}(0)\).

Problem (11B.17)

Suppose \(p \in [1, \infty]\) and \(f \in L^p(\partial D)\). Prove that if \(0 \le r < s < 1\), then

\begin{equation*} \|\mathcal{P}_r f\|_p \le \|\mathcal{P}_s f\|_p . \end{equation*}

Solution

Because \(\sigma\) is a finite measure with \(\sigma(\partial D) = 1\), Hölder’s inequality (7.9) gives \(\|u\|_1 \le \|u\|_p\) for every \(u \in L^p(\partial D)\); hence \(L^p(\partial D) \subseteq L^1(\partial D)\) and in particular \(f \in L^1(\partial D)\), so \(\mathcal{P}_t f\) is defined for every \(t \in [0,1)\).

Step 1: \(\mathcal{P}_t\) does not increase the \(L^p\)-norm. We claim that

\begin{equation*} \|\mathcal{P}_t g\|_p \le \|g\|_p \qquad \text{for all } t \in [0,1) \text{ and all } g \in L^p(\partial D). \end{equation*}

Indeed, 11.35 is the integral formula \((\mathcal{P}_t g)(z) = \int_{\partial D} g(w) P_t(z \bar{w}) \, d\sigma(w)\), which in the convolution notation of 11.36 says exactly that \(\mathcal{P}_t g = g * P_t\); and convolution is commutative (11.41), so \(\mathcal{P}_t g = P_t * g\). The Poisson kernel \(P_t\) is a continuous, strictly positive function on \(\partial D\), hence lies in \(L^1(\partial D)\), with

\begin{equation*} \|P_t\|_1 = \int_{\partial D} |P_t| \, d\sigma = \int_{\partial D} P_t \, d\sigma = 1, \end{equation*}

where the middle equality holds by 11.16(a) and the last equality is 11.16(b). Applying 11.38 with the \(L^1\)-function \(P_t\) and the \(L^p\)-function \(g\) gives

\begin{equation*} \|\mathcal{P}_t g\|_p = \|P_t * g\|_p \le \|P_t\|_1 \, \|g\|_p = \|g\|_p , \end{equation*}

which proves the claim.

Step 2: factoring \(\mathcal{P}_r\) through \(\mathcal{P}_s\). Suppose \(0 \le r < s < 1\). Then \(s > 0\) (because \(s > r \ge 0\)), so we may set

\begin{equation*} t = \frac{r}{s} . \end{equation*}

Since \(0 \le r < s\), we have \(t \in [0,1)\). By the previous exercise (11B.16), applied to \(f \in L^1(\partial D)\) with \(t\) and \(s\) in the roles of \(r\) and \(s\),

\begin{equation*} \mathcal{P}_t(\mathcal{P}_s f) = \mathcal{P}_{ts} f = \mathcal{P}_r f , \end{equation*}

because \(ts = (r/s)s = r\).

Step 3: conclusion. The function \(\mathcal{P}_s f\) is continuous on \(\partial D\) (as noted in the discussion accompanying the definition 11.11, the series defining it converges uniformly on \(\partial D\)), hence bounded; because \(\sigma\) is a finite measure, this implies \(\mathcal{P}_s f \in L^p(\partial D)\) for every \(p \in [1,\infty]\). Therefore Step 1 may be applied with \(g = \mathcal{P}_s f\), and combining it with Step 2 yields

\begin{equation*} \|\mathcal{P}_r f\|_p = \bigl\| \mathcal{P}_t(\mathcal{P}_s f) \bigr\|_p \le \| \mathcal{P}_s f \|_p , \end{equation*}

as desired.

[Interpretation: for each \(\rho \in [0,1)\) the function \(\mathcal{P}_\rho f\) on \(\partial D\) records the values of the Poisson integral of \(f\) on the circle of radius \(\rho\) inside \(D\). The exercise says that \(\rho \mapsto \|\mathcal{P}_\rho f\|_p\) is nondecreasing on \([0,1)\). Step 1 bounds it above by \(\|f\|_p\), and for \(p \in [1,\infty)\) result 11.42 gives \(\|f - \mathcal{P}_\rho f\|_p \to 0\) as \(\rho \uparrow 1\), hence \(\|\mathcal{P}_\rho f\|_p \to \|f\|_p\); so in that range these norms increase to \(\|f\|_p\).]

Problem (11B.18)

Prove Wirtinger’s inequality: If \(f \colon \mathbb{R} \to \mathbb{R}\) is a continuously differentiable \(2\pi\)-periodic function and \(\int_{-\pi}^{\pi} f(t) \, dt = 0\), then

\begin{equation*} \int_{-\pi}^{\pi} \bigl( f(t) \bigr)^2 \, dt \le \int_{-\pi}^{\pi} \bigl( f^{\prime}(t) \bigr)^2 \, dt, \end{equation*}

with equality if and only if \(f(t) = a \sin(t) + b \cos(t)\) for some constants \(a, b\).

Solution

Transfer \(f\) to the unit circle. Because \(t \mapsto e^{it}\) maps \((-\pi, \pi]\) one-to-one onto \(\partial D\) and \(f\) is \(2\pi\)-periodic, there is a well-defined function \(F \colon \partial D \to \mathbb{R}\) given by

\begin{equation*} F(e^{it}) = f(t) \qquad (t \in \mathbb{R}). \end{equation*}

In the notation of 11.24 we then have \(\widetilde{F} = f\). Since \(f\) is continuously differentiable on \(\mathbb{R}\), the function \(F\) is once continuously differentiable in the sense of 11.24, and \(F^{[1]} \colon \partial D \to \mathbb{R}\) is given by \(F^{[1]}(e^{it}) = f^{\prime}(t)\).

Both \(F\) and \(F^{[1]}\) are continuous on \(\partial D\), hence bounded; because \(\sigma(\partial D) = 1\), both belong to \(L^2(\partial D)\). Recall also that for a measurable \(u \colon \partial D \to \mathbb{C}\),

\begin{equation*} \int_{\partial D} u \, d\sigma = \int_{-\pi}^{\pi} u(e^{it}) \, \frac{dt}{2\pi} . \end{equation*}

In particular

\begin{equation*} \|F\|_2^2 = \frac{1}{2\pi} \int_{-\pi}^{\pi} \bigl( f(t) \bigr)^2 \, dt, \qquad \bigl\| F^{[1]} \bigr\|_2^2 = \frac{1}{2\pi} \int_{-\pi}^{\pi} \bigl( f^{\prime}(t) \bigr)^2 \, dt , \end{equation*}

so the desired inequality is exactly \(\|F\|_2^2 \le \|F^{[1]}\|_2^2\).

The two hypotheses in Fourier terms. First, the mean-zero hypothesis says precisely that

\begin{equation*} \hat{F}(0) = \int_{\partial D} F \, d\sigma = \frac{1}{2\pi} \int_{-\pi}^{\pi} f(t) \, dt = 0 . \end{equation*}

Second, 11.26 with \(k = 1\) gives

\begin{equation*} \widehat{F^{[1]}}(n) = i n \hat{F}(n) \qquad \text{for every } n \in \mathbb{Z}, \end{equation*}

so that \(\bigl| \widehat{F^{[1]}}(n) \bigr|^2 = n^2 \bigl| \hat{F}(n) \bigr|^2\). (This is just integration by parts: for \(n \in \mathbb{Z}\), \(\int_{-\pi}^{\pi} f^{\prime}(t) e^{-int} \frac{dt}{2\pi} = \bigl[ f(t) e^{-int} \bigr]_{t=-\pi}^{t=\pi} \frac{1}{2\pi} + in \int_{-\pi}^{\pi} f(t) e^{-int} \frac{dt}{2\pi}\), and the boundary term vanishes because \(f\) and \(t \mapsto e^{-int}\) are both \(2\pi\)-periodic.)

Proof of the inequality. By 11.30 the family \(\{z^n\}_{n \in \mathbb{Z}}\) is an orthonormal basis of \(L^2(\partial D)\), so Parseval’s identity [8.63(c)] applied to \(F\) and to \(F^{[1]}\) gives

\begin{equation*} \|F\|_2^2 = \sum_{n=-\infty}^{\infty} \bigl| \hat{F}(n) \bigr|^2 = \sum_{n \ne 0} \bigl| \hat{F}(n) \bigr|^2 , \qquad \bigl\| F^{[1]} \bigr\|_2^2 = \sum_{n=-\infty}^{\infty} n^2 \bigl| \hat{F}(n) \bigr|^2 = \sum_{n \ne 0} n^2 \bigl| \hat{F}(n) \bigr|^2 , \end{equation*}

where the second equality in the first display uses \(\hat{F}(0) = 0\). Both sums are finite (each is the square of an \(L^2\)-norm of a continuous function). Subtracting,

\begin{equation*} \bigl\| F^{[1]} \bigr\|_2^2 - \|F\|_2^2 = \sum_{n \ne 0} (n^2 - 1) \bigl| \hat{F}(n) \bigr|^2 \ge 0 , \end{equation*}

because \(n^2 - 1 \ge 0\) for every nonzero integer \(n\). Multiplying by \(2\pi\) yields

\begin{equation*} \int_{-\pi}^{\pi} \bigl( f(t) \bigr)^2 \, dt \le \int_{-\pi}^{\pi} \bigl( f^{\prime}(t) \bigr)^2 \, dt , \end{equation*}

which is Wirtinger’s inequality.

The equality case, necessity. Suppose equality holds. Then \(\sum_{n \ne 0} (n^2 - 1) |\hat{F}(n)|^2 = 0\). Every term of this series is nonnegative, so every term is \(0\); since \(n^2 - 1 > 0\) whenever \(|n| \ge 2\), we conclude that

\begin{equation*} \hat{F}(n) = 0 \qquad \text{for all } n \in \mathbb{Z} \text{ with } |n| \ge 2 , \end{equation*}

and we already know \(\hat{F}(0) = 0\). Thus at most the coefficients \(\hat{F}(1)\) and \(\hat{F}(-1)\) are nonzero. By 11.31 the Fourier series of \(F\) converges to \(F\) in the norm of \(L^2(\partial D)\); as all but two of its terms vanish, this says

\begin{equation*} F(z) = \hat{F}(1) z + \hat{F}(-1) \bar{z} \qquad \text{for almost every } z \in \partial D . \end{equation*}

Both sides of the last equation are continuous on \(\partial D\). We claim that this upgrades the almost-everywhere equality to an everywhere equality. Indeed, let \(E\) be the set where the two sides differ; then \(E\) is relatively open in \(\partial D\), so \(A = \{t \in (-\pi, \pi] : e^{it} \in E\}\) is relatively open in \((-\pi, \pi]\) (the map \(t \mapsto e^{it}\) is continuous). If \(E\) were nonempty then \(A\) would be nonempty and relatively open, hence would contain an interval of positive length, so \(|A| > 0\); by the definition 11.4 of \(\sigma\) this gives \(\sigma(E) = |A| / (2\pi) > 0\), contradicting \(\sigma(E) = 0\). Hence \(E = \varnothing\) and the equation above holds for every \(z \in \partial D\). Equivalently,

\begin{equation*} f(t) = \hat{F}(1) e^{it} + \hat{F}(-1) e^{-it} \qquad \text{for all } t \in \mathbb{R} . \end{equation*}

We now use that \(f\) is real valued to relate the two surviving coefficients:

\begin{equation*} \hat{F}(-1) = \int_{-\pi}^{\pi} f(t) e^{it} \, \frac{dt}{2\pi} = \overline{ \int_{-\pi}^{\pi} f(t) e^{-it} \, \frac{dt}{2\pi} } = \overline{\hat{F}(1)} \end{equation*}

where the first equality comes from \(\hat{F}(-1) = \int_{\partial D} F(z) \bar{z}^{\,-1} \, d\sigma(z) = \int_{\partial D} F(z)\, z \, d\sigma(z)\) (using \(\bar{z}^{\,-1} = z\) on \(\partial D\)), and the second holds because \(f\) is real valued, so that \(\overline{f(t) e^{-it}} = f(t) e^{it}\). Writing \(\hat{F}(1) = \alpha - i\beta\) with \(\alpha, \beta \in \mathbb{R}\), we therefore get

\begin{equation*} f(t) = 2 \operatorname{Re}\bigl( \hat{F}(1) e^{it} \bigr) = 2 \operatorname{Re}\bigl( (\alpha - i\beta)(\cos t + i \sin t) \bigr) = 2\alpha \cos t + 2\beta \sin t . \end{equation*}

So \(f(t) = a \sin t + b \cos t\) with \(a = 2\beta\) and \(b = 2\alpha\), as claimed.

The equality case, sufficiency. Conversely, suppose \(f(t) = a \sin t + b \cos t\) for some constants \(a, b\). Such an \(f\) is indeed continuously differentiable, \(2\pi\)-periodic, and satisfies \(\int_{-\pi}^{\pi} f(t) \, dt = 0\). Using \(\int_{-\pi}^{\pi} \sin^2 t \, dt = \int_{-\pi}^{\pi} \cos^2 t \, dt = \pi\) and \(\int_{-\pi}^{\pi} \sin t \cos t \, dt = 0\), we get

\begin{equation*} \int_{-\pi}^{\pi} \bigl( f(t) \bigr)^2 \, dt = \pi ( a^2 + b^2 ) . \end{equation*}

Since \(f^{\prime}(t) = a \cos t - b \sin t\), the same three integrals give

\begin{equation*} \int_{-\pi}^{\pi} \bigl( f^{\prime}(t) \bigr)^2 \, dt = \pi ( a^2 + b^2 ) . \end{equation*}

Thus equality holds in Wirtinger’s inequality for these functions, completing the proof.

Exercises 11C

Problem (11C.1)

Suppose \(f \in L^1(\mathbb{R})\). Prove that \(\|\hat f\|_\infty = \|f\|_1\) if and only if there exists \(\zeta \in \partial\mathbb{D}\) and \(t \in \mathbb{R}\) such that \(\zeta f(x) e^{-itx} \ge 0\) for almost every \(x \in \mathbb{R}\).

Solution

Throughout, \(\hat f(s) = \int_{-\infty}^{\infty} f(x) e^{-2\pi i s x}\,dx\), and by the Riemann–Lebesgue Lemma (11.49) the function \(\hat f\) is uniformly continuous on \(\mathbb{R}\), satisfies \(\|\hat f\|_\infty \le \|f\|_1\), and has limit \(0\) at \(\pm\infty\). Because \(\hat f\) is continuous, its essential supremum coincides with its ordinary supremum:

\begin{equation*} \|\hat f\|_\infty = \sup_{s \in \mathbb{R}} |\hat f(s)| . \end{equation*}

Indeed, \(\|\hat f\|_\infty \le \sup_{s} |\hat f(s)|\) always holds; conversely, if \(|\hat f(s_1)| > c\) for some \(s_1 \in \mathbb{R}\), then by continuity \(|\hat f| > c\) on some open interval containing \(s_1\), and that interval has positive Lebesgue measure, so \(\|\hat f\|_\infty \ge c\). Letting \(c\) increase to \(\sup_s |\hat f(s)|\) gives the reverse inequality.

First suppose there exist \(\zeta \in \partial\mathbb{D}\) and \(t \in \mathbb{R}\) with \(\zeta f(x) e^{-itx} \ge 0\) for almost every \(x \in \mathbb{R}\). Put \(s_0 = t/(2\pi)\), so that \(e^{-2\pi i s_0 x} = e^{-itx}\) for all \(x\). Because \(|\zeta| = 1\) and \(|e^{-itx}| = 1\), we have \(|\zeta f(x) e^{-itx}| = |f(x)|\) for all \(x\); combined with the hypothesis that \(\zeta f(x) e^{-itx}\) is a nonnegative real number for almost every \(x\), this gives

\begin{equation*} \zeta f(x) e^{-itx} = |f(x)| \quad \text{for almost every } x \in \mathbb{R}. \end{equation*}

Hence

\begin{equation*} |\hat f(s_0)| = \Bigl| \int_{-\infty}^{\infty} f(x) e^{-itx}\,dx \Bigr| = \Bigl| \bar\zeta \int_{-\infty}^{\infty} \zeta f(x) e^{-itx}\,dx \Bigr| = \int_{-\infty}^{\infty} |f(x)|\,dx = \|f\|_1 , \end{equation*}

where we used \(|\bar\zeta| = 1\). Thus \(\|\hat f\|_\infty \ge \|f\|_1\), and the reverse inequality holds by 11.49. Therefore \(\|\hat f\|_\infty = \|f\|_1\).

Now suppose \(\|\hat f\|_\infty = \|f\|_1\). If \(\|f\|_1 = 0\) then \(f(x) = 0\) for almost every \(x\), and the desired conclusion holds with \(\zeta = 1\) and \(t = 0\) (because then \(\zeta f(x) e^{-itx} = 0 \ge 0\) for almost every \(x\)). So assume \(\|f\|_1 > 0\).

We claim the supremum defining \(\|\hat f\|_\infty\) is attained. Because \(\lim_{s \to \pm\infty} \hat f(s) = 0\) (by 11.49), there exists \(M > 0\) such that \(|\hat f(s)| < \tfrac{1}{2}\|f\|_1\) for all \(s\) with \(|s| > M\). Since \(\|\hat f\|_\infty = \|f\|_1 > \tfrac12 \|f\|_1\), we conclude that

\begin{equation*} \sup_{s \in \mathbb{R}} |\hat f(s)| = \sup_{|s| \le M} |\hat f(s)| . \end{equation*}

The function \(|\hat f|\) is continuous on the compact interval \([-M, M]\), so it attains its maximum there, at some \(s_0 \in [-M, M]\). Thus \(|\hat f(s_0)| = \|f\|_1\).

Write \(\hat f(s_0) = \lambda \|f\|_1\) where \(\lambda \in \mathbb{C}\) with \(|\lambda| = 1\) (this is possible because \(\|f\|_1 > 0\)). Set \(\zeta = \bar\lambda \in \partial\mathbb{D}\) and \(t = 2\pi s_0 \in \mathbb{R}\), and define \(g : \mathbb{R} \to \mathbb{C}\) by

\begin{equation*} g(x) = \zeta f(x) e^{-itx} . \end{equation*}

Then \(|g(x)| = |f(x)|\) for all \(x\), and

\begin{equation*} \int_{-\infty}^{\infty} g(x)\,dx = \bar\lambda \int_{-\infty}^{\infty} f(x) e^{-2\pi i s_0 x}\,dx = \bar\lambda\, \hat f(s_0) = \bar\lambda \lambda \|f\|_1 = \|f\|_1 = \int_{-\infty}^{\infty} |g(x)|\,dx . \end{equation*}

Taking imaginary parts of the two ends of the displayed chain (both integrals are finite, since \(g \in L^1(\mathbb{R})\)) gives \(\int_{-\infty}^{\infty} \operatorname{Im} g(x)\,dx = 0\), and taking real parts gives

\begin{equation*} \int_{-\infty}^{\infty} \bigl( |g(x)| - \operatorname{Re} g(x) \bigr)\,dx = 0 . \end{equation*}

Write \(h = |g| - \operatorname{Re} g\). Then \(h\) is a nonnegative function in \(L^1(\mathbb{R})\) (nonnegative because \(\operatorname{Re} w \le |w|\) for every \(w \in \mathbb{C}\)) with \(\|h\|_1 = \int h = 0\). Markov’s inequality (4.1) gives \(|\{x \in \mathbb{R} : h(x) \ge c\}| \le \frac1c \|h\|_1 = 0\) for every \(c > 0\); taking \(c = 1/n\) and letting \(n \to \infty\) (using countable subadditivity of Lebesgue measure) shows \(|\{x : h(x) > 0\}| = 0\). Thus \(h = 0\) almost everywhere, that is,

\begin{equation*} \operatorname{Re} g(x) = |g(x)| \quad \text{for almost every } x \in \mathbb{R}. \end{equation*}

For each such \(x\) we have \((\operatorname{Re} g(x))^2 + (\operatorname{Im} g(x))^2 = |g(x)|^2 = (\operatorname{Re} g(x))^2\), so \(\operatorname{Im} g(x) = 0\), and therefore \(g(x) = \operatorname{Re} g(x) = |g(x)| \ge 0\).

In other words, \(\zeta f(x) e^{-itx} \ge 0\) for almost every \(x \in \mathbb{R}\), with \(\zeta \in \partial\mathbb{D}\) and \(t \in \mathbb{R}\) as constructed. This completes the proof.

Problem (11C.2)

Suppose \(f(x) = x e^{-\pi x^2}\) for all \(x \in \mathbb{R}\). Show that \(\hat f = -i f\).

Solution

Define \(\varphi : \mathbb{R} \to \mathbb{R}\) by \(\varphi(x) = e^{-\pi x^2}\). Then \(\varphi \in L^1(\mathbb{R})\) and \(f(x) = x \varphi(x)\), so \(f \in L^1(\mathbb{R})\) as well (because \(|x| e^{-\pi x^2}\) is continuous and decays faster than any power of \(1/|x|\); explicitly \(\int_{-\infty}^{\infty} |x| e^{-\pi x^2}\,dx = 2\int_0^\infty x e^{-\pi x^2}\,dx = \frac1\pi < \infty\)).

Example 11.51 shows that \(\hat\varphi = \varphi\).

Now apply 11.50 with the function \(\varphi\) in the role of \(f\) there and with \(g(x) = x\varphi(x) = f(x)\). Since \(g = f \in L^1(\mathbb{R})\), the hypothesis of 11.50 is satisfied, and 11.50 tells us that \(\hat\varphi\) is continuously differentiable on \(\mathbb{R}\) with

\begin{equation*} (\hat\varphi)^{\prime}(t) = -2\pi i\, \hat f(t) \end{equation*}

for all \(t \in \mathbb{R}\).

Because \(\hat\varphi = \varphi\), the left side equals

\begin{equation*} \varphi^{\prime}(t) = \frac{d}{dt} e^{-\pi t^2} = -2\pi t\, e^{-\pi t^2}. \end{equation*}

Therefore \(-2\pi t e^{-\pi t^2} = -2\pi i \hat f(t)\) for all \(t \in \mathbb{R}\), and dividing by \(-2\pi i\) gives

\begin{equation*} \hat f(t) = \frac{-2\pi t e^{-\pi t^2}}{-2\pi i} = \frac{1}{i}\, t e^{-\pi t^2} = -i\, t e^{-\pi t^2} = -i f(t) \end{equation*}

for all \(t \in \mathbb{R}\), where we used \(1/i = -i\). Thus \(\hat f = -i f\), as desired.

In particular \(f\) is an eigenvector of the Fourier transform with eigenvalue \(-i\), which is the fact used in the proof of 11.87(c).

Problem (11C.3)

Suppose \(f(x) = 4\pi x^2 e^{-\pi x^2} - e^{-\pi x^2}\) for all \(x \in \mathbb{R}\). Show that \(\hat f = -f\).

Solution

Define \(\varphi, f_1, f_2 : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} \varphi(x) = e^{-\pi x^2}, \qquad f_1(x) = x e^{-\pi x^2}, \qquad f_2(x) = x^2 e^{-\pi x^2} . \end{equation*}

Each of these functions is continuous and decays faster than any power of \(1/|x|\), so all three are in \(L^1(\mathbb{R})\). Note that

\begin{equation*} f = 4\pi f_2 - \varphi . \end{equation*}

Because the Fourier transform is linear, it suffices to compute \(\hat\varphi\) and \(\hat f_2\).

By Example 11.51 we have \(\hat\varphi = \varphi\), and by Exercise 11C.2 we have

\begin{equation*} \hat f_1(t) = -i\, t e^{-\pi t^2} \end{equation*}

for all \(t \in \mathbb{R}\).

Now apply 11.50 with \(f_1\) in the role of \(f\) there and with \(g(x) = x f_1(x) = x^2 e^{-\pi x^2} = f_2(x)\). Since \(f_2 \in L^1(\mathbb{R})\), result 11.50 gives

\begin{equation*} (\hat f_1)^{\prime}(t) = -2\pi i\, \hat f_2(t) \end{equation*}

for all \(t \in \mathbb{R}\). Differentiating the formula for \(\hat f_1\) gives

\begin{equation*} (\hat f_1)^{\prime}(t) = \frac{d}{dt}\bigl( -i t e^{-\pi t^2} \bigr) = -i e^{-\pi t^2} - i t \cdot (-2\pi t) e^{-\pi t^2} = \bigl( -i + 2\pi i t^2 \bigr) e^{-\pi t^2} . \end{equation*}

Hence

\begin{equation*} \hat f_2(t) = \frac{(-i + 2\pi i t^2) e^{-\pi t^2}}{-2\pi i} = \frac{-1 + 2\pi t^2}{-2\pi}\, e^{-\pi t^2} = \Bigl( \frac{1}{2\pi} - t^2 \Bigr) e^{-\pi t^2} \end{equation*}

for all \(t \in \mathbb{R}\).

Therefore, for all \(t \in \mathbb{R}\),

\begin{equation*} \begin{aligned} \hat f(t) &= 4\pi\, \hat f_2(t) - \hat\varphi(t) \\ &= 4\pi \Bigl( \frac{1}{2\pi} - t^2 \Bigr) e^{-\pi t^2} - e^{-\pi t^2} \\ &= \bigl( 2 - 4\pi t^2 \bigr) e^{-\pi t^2} - e^{-\pi t^2} \\ &= \bigl( 1 - 4\pi t^2 \bigr) e^{-\pi t^2} \\ &= -\bigl( 4\pi t^2 e^{-\pi t^2} - e^{-\pi t^2} \bigr) \\ &= -f(t) . \end{aligned} \end{equation*}

Thus \(\hat f = -f\), so \(f\) is an eigenvector of the Fourier transform with eigenvalue \(-1\).

Problem (11C.4)

Find \(f \in L^1(\mathbb{R})\) such that \(f \ne 0\) and \(\hat f = i f\).

Solution

Take

\begin{equation*} f(x) = \bigl( 4\pi x^3 - 3x \bigr) e^{-\pi x^2} . \end{equation*}

This function is continuous and decays faster than any power of \(1/|x|\), so \(f \in L^1(\mathbb{R})\), and clearly \(f \ne 0\) (for instance \(f(1) = (4\pi - 3)e^{-\pi} \ne 0\)). We claim \(\hat f = i f\).

Define \(f_1, f_2, f_3 : \mathbb{R} \to \mathbb{R}\) by \(f_k(x) = x^k e^{-\pi x^2}\); each is in \(L^1(\mathbb{R})\) for the same reason. Then \(f = 4\pi f_3 - 3 f_1\).

By Exercise 11C.2,

\begin{equation*} \hat f_1(t) = -i\, t e^{-\pi t^2} , \end{equation*}

and, as computed in the solution of Exercise 11C.3 (apply 11.50 to \(f_1\) with \(g = f_2\)),

\begin{equation*} \hat f_2(t) = \Bigl( \frac{1}{2\pi} - t^2 \Bigr) e^{-\pi t^2} . \end{equation*}

Now apply 11.50 with \(f_2\) in the role of \(f\) there and \(g(x) = x f_2(x) = f_3(x)\); since \(f_3 \in L^1(\mathbb{R})\), we get \((\hat f_2)^{\prime}(t) = -2\pi i\, \hat f_3(t)\) for all \(t \in \mathbb{R}\). Differentiating,

\begin{equation*} (\hat f_2)^{\prime}(t) = -2t\, e^{-\pi t^2} + \Bigl( \frac{1}{2\pi} - t^2 \Bigr)(-2\pi t) e^{-\pi t^2} = \bigl( -2t - t + 2\pi t^3 \bigr) e^{-\pi t^2} = \bigl( 2\pi t^3 - 3t \bigr) e^{-\pi t^2} . \end{equation*}

Hence, using \(\dfrac{1}{-2\pi i} = \dfrac{i}{2\pi}\),

\begin{equation*} \hat f_3(t) = \frac{(2\pi t^3 - 3t) e^{-\pi t^2}}{-2\pi i} = \frac{i}{2\pi}\bigl( 2\pi t^3 - 3t \bigr) e^{-\pi t^2} = i\Bigl( t^3 - \frac{3t}{2\pi} \Bigr) e^{-\pi t^2} \end{equation*}

for all \(t \in \mathbb{R}\).

Therefore, for all \(t \in \mathbb{R}\),

\begin{equation*} \begin{aligned} \hat f(t) &= 4\pi\, \hat f_3(t) - 3\, \hat f_1(t) \\ &= 4\pi i \Bigl( t^3 - \frac{3t}{2\pi} \Bigr) e^{-\pi t^2} - 3 \bigl( -i t e^{-\pi t^2} \bigr) \\ &= i\bigl( 4\pi t^3 - 6t \bigr) e^{-\pi t^2} + 3 i t e^{-\pi t^2} \\ &= i \bigl( 4\pi t^3 - 3t \bigr) e^{-\pi t^2} \\ &= i f(t) . \end{aligned} \end{equation*}

Thus \(\hat f = i f\) with \(f \in L^1(\mathbb{R})\) and \(f \ne 0\), as required. (This function is the third Hermite function; together with Example 11.51 and Exercises 11C.2 and 11C.3 it exhibits eigenvectors of the Fourier transform for each of the four eigenvalues \(1, -i, -1, i\), as used in the proof of 11.87(c).)

Problem (11C.5)

Prove that if \(p\) is a polynomial on \(\mathbb{R}\) with complex coefficients and \(f : \mathbb{R} \to \mathbb{C}\) is defined by \(f(x) = p(x) e^{-\pi x^2}\), then there exists a polynomial \(q\) on \(\mathbb{R}\) with complex coefficients such that \(\deg q = \deg p\) and \(\hat f(t) = q(t) e^{-\pi t^2}\) for all \(t \in \mathbb{R}\).

Solution

For each nonnegative integer \(n\) define \(g_n : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} g_n(x) = x^n e^{-\pi x^2} . \end{equation*}

Each \(g_n\) is continuous and \(|g_n(x)| \le |x|^n e^{-\pi x^2}\), which tends to \(0\) faster than any power of \(1/|x|\) as \(|x| \to \infty\); hence \(g_n \in L^1(\mathbb{R})\) for every \(n\).

Claim: for each nonnegative integer \(n\) there is a polynomial \(q_n\) with complex coefficients such that

\begin{equation*} \hat g_n(t) = q_n(t) e^{-\pi t^2} \quad \text{for all } t \in \mathbb{R}, \end{equation*}

with \(\deg q_n = n\) and with leading coefficient of \(q_n\) equal to \((-i)^n\).

We prove the claim by induction on \(n\).

For \(n = 0\) we have \(g_0(x) = e^{-\pi x^2}\), and Example 11.51 gives \(\hat g_0(t) = e^{-\pi t^2}\). So \(q_0 = 1\), which has degree \(0\) and leading coefficient \(1 = (-i)^0\).

Now suppose \(n \ge 0\) and the claim holds for \(n\); we prove it for \(n + 1\). Apply 11.50 with \(g_n\) in the role of \(f\) there and with \(g(x) = x g_n(x) = g_{n+1}(x)\). Since \(g_n \in L^1(\mathbb{R})\) and \(g_{n+1} \in L^1(\mathbb{R})\), result 11.50 says that \(\hat g_n\) is continuously differentiable on \(\mathbb{R}\) and

\begin{equation*} (\hat g_n)^{\prime}(t) = -2\pi i\, \hat g_{n+1}(t) \end{equation*}

for all \(t \in \mathbb{R}\). By the induction hypothesis,

\begin{equation*} (\hat g_n)^{\prime}(t) = \frac{d}{dt}\bigl( q_n(t) e^{-\pi t^2} \bigr) = \bigl( q_n^{\prime}(t) - 2\pi t\, q_n(t) \bigr) e^{-\pi t^2} . \end{equation*}

Hence

\begin{equation*} \hat g_{n+1}(t) = q_{n+1}(t) e^{-\pi t^2}, \qquad \text{where } q_{n+1}(t) = \frac{q_n^{\prime}(t) - 2\pi t\, q_n(t)}{-2\pi i} = \frac{i}{2\pi}\bigl( q_n^{\prime}(t) - 2\pi t\, q_n(t) \bigr) . \end{equation*}

Clearly \(q_{n+1}\) is a polynomial with complex coefficients. Now \(\deg q_n^{\prime} \le n - 1 < n + 1\) (interpret \(q_0^{\prime} = 0\) as having no degree), while \(t \mapsto t\, q_n(t)\) has degree \(n+1\) with leading coefficient \((-i)^n\). Therefore \(q_{n+1}\) has degree exactly \(n+1\), with leading coefficient

\begin{equation*} \frac{i}{2\pi} \cdot (-2\pi) \cdot (-i)^n = -i \cdot (-i)^n = (-i)^{n+1} \ne 0 . \end{equation*}

This completes the induction and proves the claim.

Now let \(p\) be a polynomial on \(\mathbb{R}\) with complex coefficients. If \(p = 0\) then \(f = 0\), so \(\hat f = 0\), and taking \(q = 0\) gives \(\hat f(t) = q(t) e^{-\pi t^2}\) for all \(t\) together with \(\deg q = \deg p\) (both sides being the degree of the zero polynomial, whatever convention is used for it). So assume \(p \ne 0\), say

\begin{equation*} p(x) = \sum_{k=0}^{n} a_k x^k \quad \text{with } a_n \ne 0, \text{ so } \deg p = n . \end{equation*}

Then \(f = \sum_{k=0}^n a_k g_k\), and each \(g_k \in L^1(\mathbb{R})\), so \(f \in L^1(\mathbb{R})\). Because the Fourier transform is linear (immediate from its definition and the linearity of the integral),

\begin{equation*} \hat f(t) = \sum_{k=0}^{n} a_k \hat g_k(t) = \Bigl( \sum_{k=0}^{n} a_k q_k(t) \Bigr) e^{-\pi t^2} \quad \text{for all } t \in \mathbb{R}. \end{equation*}

Set \(q = \sum_{k=0}^n a_k q_k\), a polynomial with complex coefficients satisfying \(\hat f(t) = q(t) e^{-\pi t^2}\) for all \(t \in \mathbb{R}\).

Finally we check \(\deg q = n\). For \(k < n\) we have \(\deg(a_k q_k) \le k \le n-1\), while \(a_n q_n\) has degree exactly \(n\) with leading coefficient \(a_n(-i)^n\), which is nonzero because \(a_n \ne 0\). Hence the coefficient of \(t^n\) in \(q\) equals \(a_n (-i)^n \ne 0\) and no higher power of \(t\) occurs, so \(\deg q = n = \deg p\), as desired.

Problem (11C.6)

Suppose

\begin{equation*} f(x) = \begin{cases} x e^{-2\pi x} & \text{if } x > 0, \\ 0 & \text{if } x \le 0. \end{cases} \end{equation*}

Show that \(\hat f(t) = \dfrac{1}{4\pi^2 (1 + it)^2}\) for all \(t \in \mathbb{R}\).

Solution

First note \(f \in L^1(\mathbb{R})\): the function \(f\) is measurable (it is continuous), and

\begin{equation*} \int_{-\infty}^{\infty} |f(x)|\,dx = \int_0^\infty x e^{-2\pi x}\,dx = \frac{1}{4\pi^2} < \infty , \end{equation*}

by the computation carried out below with \(a = 2\pi\).

We first record an elementary integral. Suppose \(a \in \mathbb{C}\) with \(\operatorname{Re} a > 0\). Define \(F : [0,\infty) \to \mathbb{C}\) by

\begin{equation*} F(x) = -\frac{x}{a} e^{-ax} - \frac{1}{a^2} e^{-ax} . \end{equation*}

Then

\begin{equation*} F^{\prime}(x) = -\frac{1}{a} e^{-ax} + x e^{-ax} + \frac{1}{a} e^{-ax} = x e^{-ax} . \end{equation*}

(Here \(F^{\prime}\) denotes the derivative of the complex-valued function \(F\); the computation is the product rule applied to the real and imaginary parts of \(F\).) Also \(|e^{-ax}| = e^{-x \operatorname{Re} a}\), so

\begin{equation*} |F(x)| \le \Bigl( \frac{x}{|a|} + \frac{1}{|a|^2} \Bigr) e^{-x \operatorname{Re} a} \longrightarrow 0 \quad \text{as } x \to \infty . \end{equation*}

Write \(c = \operatorname{Re} a > 0\). Repeating the computation of \(F^{\prime}\) above with \(a\) replaced by the positive real number \(c\), and using the Fundamental Theorem of Calculus on \([0,M]\), shows that \(\int_0^M x e^{-cx}\,dx = \bigl( \frac{M}{c} + \frac{1}{c^2} \bigr)\bigl( -e^{-cM} \bigr) + \frac{1}{c^2}\), which increases to \(1/c^2\) as \(M \to \infty\). Hence the Monotone Convergence Theorem (3.11), applied to the increasing sequence of nonnegative functions \(\chi_{(0,k)}(x)\, x e^{-cx}\), shows that \(x \mapsto |x e^{-ax}| = x e^{-cx}\) is integrable on \((0,\infty)\) with integral \(1/c^2\). Now the Fundamental Theorem of Calculus applied to the real and imaginary parts of \(F\) on \([0,M]\) gives \(\int_0^M x e^{-ax}\,dx = F(M) - F(0)\), and the Dominated Convergence Theorem (3.31), with dominating function \(x e^{-cx}\), lets us pass to the limit as \(M \to \infty\):

\begin{equation*} \int_0^\infty x e^{-ax}\,dx = \lim_{M \to \infty} \bigl( F(M) - F(0) \bigr) = 0 - \Bigl( -\frac{1}{a^2} \Bigr) = \frac{1}{a^2} . \end{equation*}

Now fix \(t \in \mathbb{R}\). Since \(f\) vanishes on \((-\infty, 0]\),

\begin{equation*} \hat f(t) = \int_{-\infty}^{\infty} f(x) e^{-2\pi i t x}\,dx = \int_0^\infty x e^{-2\pi x} e^{-2\pi i t x}\,dx = \int_0^\infty x e^{-2\pi(1 + it)x}\,dx . \end{equation*}

Apply the computation above with \(a = 2\pi(1 + it)\); note \(\operatorname{Re} a = 2\pi > 0\), so the hypothesis is satisfied. We obtain

\begin{equation*} \hat f(t) = \frac{1}{a^2} = \frac{1}{\bigl( 2\pi (1+it) \bigr)^2} = \frac{1}{4\pi^2 (1 + it)^2} , \end{equation*}

as desired.

Problem (11C.7)

Prove the formulas in 11.55 for the Fourier transforms of translations, rotations, and dilations.

Solution

Result 11.55 states: suppose \(f \in L^1(\mathbb{R})\), \(b \in \mathbb{R}\), and \(t \in \mathbb{R}\); then

(a) if \(g(x) = f(x - b)\) for all \(x \in \mathbb{R}\), then \(\hat g(t) = e^{-2\pi i b t} \hat f(t)\);

(b) if \(g(x) = e^{2\pi i b x} f(x)\) for all \(x \in \mathbb{R}\), then \(\hat g(t) = \hat f(t - b)\);

(c) if \(b \ne 0\) and \(g(x) = f(bx)\) for all \(x \in \mathbb{R}\), then \(\hat g(t) = \frac{1}{|b|} \hat f\bigl( \frac{t}{b} \bigr)\).

Recall the definition of the Fourier transform: for \(h \in L^1(\mathbb{R})\) and \(t \in \mathbb{R}\),

\begin{equation*} \hat h(t) = \int_{-\infty}^{\infty} h(x) e^{-2\pi i t x}\,dx . \end{equation*}

(a) Let \(g(x) = f(x - b)\). Outer measure is translation invariant (2.7). Translation by \(b\) carries Borel sets to Borel sets (it is a homeomorphism) and sets of outer measure \(0\) to sets of outer measure \(0\), so by the definition of Lebesgue measurability (2.70) it carries Lebesgue measurable sets to Lebesgue measurable sets; consequently \(g\) is measurable, since \(\{x : g(x) \in B\} = b + \{u : f(u) \in B\}\) for every Borel set \(B \subseteq \mathbb{C}\). Also \(\int_{-\infty}^{\infty} |g| = \int_{-\infty}^{\infty} |f| < \infty\), so \(g \in L^1(\mathbb{R})\). Moreover for every \(\psi \in L^1(\mathbb{R})\) we have \(\int_{-\infty}^{\infty} \psi(u + b)\,du = \int_{-\infty}^{\infty} \psi(x)\,dx\) (this holds for characteristic functions of measurable sets by translation invariance, hence for simple functions by linearity, hence for nonnegative measurable functions by the definition of the integral as a supremum over simple functions below them, hence for all of \(L^1(\mathbb{R})\) by splitting into real and imaginary, positive and negative, parts).

Applying this with \(\psi(x) = f(x) e^{-2\pi i t x}\), and using the change of variable \(x = u + b\), we get

\begin{equation*} \begin{aligned} \hat g(t) &= \int_{-\infty}^{\infty} f(x - b)\, e^{-2\pi i t x}\,dx \\ &= \int_{-\infty}^{\infty} f(u)\, e^{-2\pi i t (u + b)}\,du \\ &= e^{-2\pi i b t} \int_{-\infty}^{\infty} f(u)\, e^{-2\pi i t u}\,du \\ &= e^{-2\pi i b t} \hat f(t) , \end{aligned} \end{equation*}

where in the second line we used translation invariance applied to the function \(u \mapsto f(u) e^{-2\pi i t(u+b)}\), whose value at \(u\) equals \(\psi(u+b)\) for \(\psi(x) = f(x)e^{-2\pi i t x}\), and where \(e^{-2\pi i b t}\) came out of the integral because it does not depend on \(u\).

(b) Let \(g(x) = e^{2\pi i b x} f(x)\). Since \(|e^{2\pi i b x}| = 1\) for all \(x\), we have \(|g| = |f|\), so \(g\) is measurable and \(g \in L^1(\mathbb{R})\). Then

\begin{equation*} \begin{aligned} \hat g(t) &= \int_{-\infty}^{\infty} e^{2\pi i b x} f(x)\, e^{-2\pi i t x}\,dx \\ &= \int_{-\infty}^{\infty} f(x)\, e^{-2\pi i (t - b) x}\,dx \\ &= \hat f(t - b) , \end{aligned} \end{equation*}

where the second equality holds because \(e^{2\pi i b x} e^{-2\pi i t x} = e^{-2\pi i (t-b)x}\).

(c) Suppose \(b \ne 0\) and \(g(x) = f(bx)\). First we record the dilation rule for integrals: if \(\psi \in L^1(\mathbb{R})\) and \(b \ne 0\), then \(x \mapsto \psi(bx)\) is in \(L^1(\mathbb{R})\) and

\begin{equation*} \int_{-\infty}^{\infty} \psi(bx)\,dx = \frac{1}{|b|} \int_{-\infty}^{\infty} \psi(u)\,du . \end{equation*}

To prove this, we first record how Lebesgue measure behaves under dilation by an arbitrary nonzero real number \(c\) (the book’s 5.41 covers only \(c > 0\)). Outer measure is defined (2.2) as the infimum of \(\sum_k \ell(I_k)\) over countable covers of a set by open intervals. For \(c \ne 0\) the map \(I \mapsto cI\) is a bijection of the open intervals onto themselves satisfying \(\ell(cI) = |c|\,\ell(I)\), and \(E \subseteq \bigcup_k I_k\) if and only if \(cE \subseteq \bigcup_k cI_k\). Hence this map is a bijection between the covers of \(E\) and the covers of \(cE\) that multiplies every cover sum by \(|c|\), so

\begin{equation*} |cE| = |c| \, |E| \quad \text{for every } E \subseteq \mathbb{R}, \end{equation*}

where \(|\cdot|\) denotes outer measure. Furthermore \(cE\) is Lebesgue measurable whenever \(E\) is: by the definition of Lebesgue measurability (2.70) there is a Borel set \(B \subseteq E\) with \(|E \setminus B| = 0\); then \(cB \subseteq cE\) is a Borel set (it is the inverse image of \(B\) under the continuous map \(x \mapsto x/c\)) and \(|cE \setminus cB| = |c(E \setminus B)| = |c| \cdot 0 = 0\).

Consequently, if \(\psi : \mathbb{R} \to [0,\infty]\) is Lebesgue measurable then so is \(x \mapsto \psi(bx)\), because \(\{x : \psi(bx) \in B\} = \frac1b \{u : \psi(u) \in B\}\) for every Borel set \(B \subseteq [0,\infty]\). If \(\psi = \chi_E\) with \(E\) Lebesgue measurable, then \(\psi(bx) = \chi_{\frac1b E}(x)\), so

\begin{equation*} \int_{-\infty}^{\infty} \psi(bx)\,dx = \Bigl| \tfrac1b E \Bigr| = \frac{1}{|b|}\,|E| = \frac{1}{|b|} \int_{-\infty}^{\infty} \psi(u)\,du \end{equation*}

(both sides may equal \(\infty\), and the identity still holds). The formula extends to nonnegative simple measurable functions by linearity of the integral, then to every nonnegative measurable \(\psi\) by choosing an increasing sequence \(s_1 \le s_2 \le \cdots\) of nonnegative simple measurable functions converging pointwise to \(\psi\) (2.89), noting that \(x \mapsto s_k(bx)\) is then an increasing sequence of nonnegative simple measurable functions converging pointwise to \(x \mapsto \psi(bx)\), and applying the Monotone Convergence Theorem (3.11) to both sides. Finally it extends to all \(\psi \in L^1(\mathbb{R})\) by decomposing \(\psi\) into real and imaginary, positive and negative, parts.

In particular \(g \in L^1(\mathbb{R})\), since \(\int |g| = \frac{1}{|b|}\int |f| < \infty\). Now fix \(t \in \mathbb{R}\) and apply the dilation rule to \(\psi(u) = f(u) e^{-2\pi i (t/b) u}\), noting that \(\psi(bx) = f(bx) e^{-2\pi i t x}\):

\begin{equation*} \hat g(t) = \int_{-\infty}^{\infty} f(bx)\, e^{-2\pi i t x}\,dx = \int_{-\infty}^{\infty} \psi(bx)\,dx = \frac{1}{|b|} \int_{-\infty}^{\infty} f(u)\, e^{-2\pi i (t/b) u}\,du = \frac{1}{|b|}\, \hat f\Bigl( \frac{t}{b} \Bigr) . \end{equation*}

This proves (a), (b), and (c), which are exactly the formulas asserted in 11.55.

Problem (11C.8)

Suppose \(f \in L^1(\mathbb{R})\) and \(n \in \mathbb{Z}^+\). Define \(g \colon \mathbb{R} \to \mathbb{C}\) by \(g(x) = x^n f(x)\). Prove that if \(g \in L^1(\mathbb{R})\), then \(\hat{f}\) is \(n\) times continuously differentiable on \(\mathbb{R}\) and

\begin{equation*} (\hat{f})^{(n)}(t) = (-2\pi i)^n \hat{g}(t) \end{equation*}

for all \(t \in \mathbb{R}\).

Solution

Assume \(g \in L^1(\mathbb{R})\), where \(g(x) = x^n f(x)\).

For \(k \in \{0, 1, \dots, n\}\) define \(g_k \colon \mathbb{R} \to \mathbb{C}\) by

\begin{equation*} g_k(x) = x^k f(x) . \end{equation*}

Thus \(g_0 = f\) and \(g_n = g\).

Step 1: each \(g_k\) is in \(L^1(\mathbb{R})\). If \(0 \le k \le n\) and \(x \in \mathbb{R}\), then \(|x|^k \le 1\) when \(|x| \le 1\) and \(|x|^k \le |x|^n\) when \(|x| \ge 1\); hence \(|x|^k \le 1 + |x|^n\) for every \(x \in \mathbb{R}\). Therefore

\begin{equation*} |g_k(x)| = |x|^k \, |f(x)| \le |f(x)| + |x|^n |f(x)| = |f(x)| + |g(x)| \end{equation*}

for all \(x \in \mathbb{R}\). Because \(f\) and \(g\) are in \(L^1(\mathbb{R})\), the right side of the inequality above is an integrable function, and each \(g_k\) is measurable; hence \(g_k \in L^1(\mathbb{R})\).

Step 2: induction. We claim that for each \(k \in \{0, 1, \dots, n\}\) the function \(\hat{f}\) is \(k\) times continuously differentiable on \(\mathbb{R}\) and

\begin{equation*} (\hat{f})^{(k)}(t) = (-2\pi i)^k \, \widehat{g_k}(t) \end{equation*}

for all \(t \in \mathbb{R}\).

If \(k = 0\), then the displayed equation reads \(\hat{f} = \widehat{g_0}\), which is true because \(g_0 = f\); the continuity of \(\hat{f}\) is part of 11.49.

Now suppose \(0 \le k < n\) and the claim holds for \(k\). The function \(g_k\) is in \(L^1(\mathbb{R})\) by Step 1, and the function \(x \mapsto x \, g_k(x)\) equals \(g_{k+1}\), which is in \(L^1(\mathbb{R})\) by Step 1. Hence 11.50, applied with \(g_k\) in place of \(f\), shows that \(\widehat{g_k}\) is a continuously differentiable function on \(\mathbb{R}\) with

\begin{equation*} (\widehat{g_k})^{\prime}(t) = -2\pi i \, \widehat{g_{k+1}}(t) \end{equation*}

for all \(t \in \mathbb{R}\). Combining this with the induction hypothesis, we see that \((\hat{f})^{(k)} = (-2\pi i)^k \widehat{g_k}\) is differentiable and

\begin{equation*} (\hat{f})^{(k+1)}(t) = (-2\pi i)^k (\widehat{g_k})^{\prime}(t) = (-2\pi i)^{k+1} \widehat{g_{k+1}}(t) \end{equation*}

for all \(t \in \mathbb{R}\). Because \(g_{k+1} \in L^1(\mathbb{R})\), the function \(\widehat{g_{k+1}}\) is continuous on \(\mathbb{R}\) (indeed uniformly continuous, by 11.49). Thus \((\hat{f})^{(k+1)}\) is continuous, which means that \(\hat{f}\) is \(k+1\) times continuously differentiable on \(\mathbb{R}\). This completes the induction.

Taking \(k = n\) in the claim shows that \(\hat{f}\) is \(n\) times continuously differentiable on \(\mathbb{R}\) and that

\begin{equation*} (\hat{f})^{(n)}(t) = (-2\pi i)^n \, \widehat{g_n}(t) = (-2\pi i)^n \, \hat{g}(t) \end{equation*}

for all \(t \in \mathbb{R}\), as desired.

Problem (11C.9)

Suppose \(n \in \mathbb{Z}^+\) and \(f \in L^1(\mathbb{R})\) is \(n\) times continuously differentiable and \(f^{(k)} \in L^1(\mathbb{R})\) for \(k = 1, \dots, n\). Prove that if \(t \in \mathbb{R}\), then

\begin{equation*} \widehat{f^{(n)}}(t) = (2\pi i t)^n \hat{f}(t) . \end{equation*}

Solution

We prove by induction on \(k\) that

\begin{equation*} \widehat{f^{(k)}}(t) = (2\pi i t)^k \hat{f}(t) \end{equation*}

for all \(t \in \mathbb{R}\) and all \(k \in \{0, 1, \dots, n\}\); the case \(k = n\) is the desired conclusion.

Here \(f^{(0)}\) means \(f\), so the case \(k = 0\) is the trivially true statement \(\hat{f}(t) = \hat{f}(t)\).

Now suppose \(0 \le k < n\) and

\begin{equation*} \widehat{f^{(k)}}(t) = (2\pi i t)^k \hat{f}(t) \end{equation*}

for all \(t \in \mathbb{R}\). We verify the hypotheses of 11.54 for the function \(f^{(k)}\).

First, \(f^{(k)} \in L^1(\mathbb{R})\): if \(k = 0\) this is the hypothesis \(f \in L^1(\mathbb{R})\), and if \(k \ge 1\) this is the hypothesis \(f^{(k)} \in L^1(\mathbb{R})\).

Second, \(f^{(k)}\) is a continuously differentiable function on \(\mathbb{R}\): because \(f\) is \(n\) times continuously differentiable and \(k + 1 \le n\), the function \(f^{(k)}\) is differentiable on \(\mathbb{R}\) with derivative \(f^{(k+1)}\), and \(f^{(k+1)}\) is continuous on \(\mathbb{R}\) (it is continuous because \(f^{(n)}\) is continuous and, if \(k+1 < n\), because \(f^{(k+1)}\) is differentiable and hence continuous).

Third, \((f^{(k)})^{\prime} = f^{(k+1)} \in L^1(\mathbb{R})\), again by hypothesis (note \(1 \le k+1 \le n\)).

Thus 11.54, applied with \(f^{(k)}\) in place of \(f\), gives

\begin{equation*} \widehat{f^{(k+1)}}(t) = \widehat{\bigl( f^{(k)} \bigr)^{\prime}}(t) = 2\pi i t \, \widehat{f^{(k)}}(t) = 2\pi i t \, (2\pi i t)^k \hat{f}(t) = (2\pi i t)^{k+1} \hat{f}(t) \end{equation*}

for all \(t \in \mathbb{R}\), where the third equality uses the induction hypothesis. This completes the induction.

Taking \(k = n\) gives \(\widehat{f^{(n)}}(t) = (2\pi i t)^n \hat{f}(t)\) for all \(t \in \mathbb{R}\), as desired.

Problem (11C.10)

Suppose \(1 \le p \le \infty\), \(f \in L^p(\mathbb{R})\), and \(g \in L^{p^{\prime}}(\mathbb{R})\). Prove that \(f * g\) is a uniformly continuous function on \(\mathbb{R}\).

Solution

Here \(p^{\prime}\) denotes the exponent conjugate to \(p\), so that \(\frac{1}{p} + \frac{1}{p^{\prime}} = 1\) (with \(1^{\prime} = \infty\) and \(\infty^{\prime} = 1\)).

Notation: for \(h \in \mathbb{R}\) and a function \(u\) on \(\mathbb{R}\), let \(u_h\) denote the translate of \(u\) defined by \(u_h(x) = u(x - h)\), as in Exercise 23 in Section 7A. Lebesgue measure is invariant under the map \(t \mapsto c - t\) for each fixed \(c \in \mathbb{R}\) (translation invariance of Lebesgue outer measure is 2.7, and invariance under \(t \mapsto -t\) is immediate from the definition of outer measure in terms of lengths of intervals, because \(\ell(-I) = \ell(I)\) for every open interval \(I\); both maps therefore also preserve the collection of Lebesgue measurable sets). Hence

\begin{equation*} \int_{-\infty}^{\infty} |u(c - t)|^q \, dt = \int_{-\infty}^{\infty} |u(s)|^q \, ds \end{equation*}

for all \(c \in \mathbb{R}\) and \(q \in [1, \infty)\), and similarly the essential supremum of \(t \mapsto |u(c-t)|\) equals \(\|u\|_\infty\). In other words, \(\|u(c - \cdot)\|_q = \|u\|_q\) for all \(q \in [1,\infty]\).

Step 1: \((f * g)(x)\) is defined for every \(x \in \mathbb{R}\). Fix \(x \in \mathbb{R}\). The functions \(t \mapsto f(t)\) and \(t \mapsto g(x - t)\) are measurable, and by the invariance noted above we have \(\|g(x - \cdot)\|_{p^{\prime}} = \|g\|_{p^{\prime}}\). Hence Hölder’s inequality (7.9) implies that \(t \mapsto f(t) g(x-t)\) is in \(L^1(\mathbb{R})\) with

\begin{equation*} \int_{-\infty}^{\infty} |f(t)\, g(x-t)| \, dt \le \|f\|_p \, \|g\|_{p^{\prime}} . \end{equation*}

Thus \((f*g)(x)\) is defined for every \(x \in \mathbb{R}\) and \(\|f*g\|_\infty \le \|f\|_p \|g\|_{p^{\prime}}\).

Step 2: reduction to a translation estimate. Because \(\frac{1}{p} + \frac{1}{p^{\prime}} = 1\), at most one of \(p\) and \(p^{\prime}\) equals \(\infty\); hence at least one of \(p, p^{\prime}\) lies in \([1, \infty)\). By 11.65 we have \(f * g = g * f\) at every point where either is defined, so by interchanging the roles of \((f,p)\) and \((g,p^{\prime})\) if necessary, we may assume that \(p < \infty\).

Now let \(x, x^{\prime} \in \mathbb{R}\) and set \(h = x^{\prime} - x\). Using 11.65 and then the substitution described above,

\begin{equation*} (f*g)(x) - (f*g)(x^{\prime}) = (g*f)(x) - (g*f)(x^{\prime}) = \int_{-\infty}^{\infty} g(t) \bigl[ f(x - t) - f(x^{\prime} - t) \bigr] \, dt . \end{equation*}

For each \(t\) we have \(f(x^{\prime} - t) = f\bigl((x-t) + h\bigr)\), so

\begin{equation*} f(x-t) - f(x^{\prime}-t) = f(x-t) - f_{-h}(x - t) = (f - f_{-h})(x-t) . \end{equation*}

Therefore Hölder’s inequality (7.9) and the invariance of \(\|\cdot\|_p\) under \(t \mapsto x - t\) give

\begin{equation*} \bigl| (f*g)(x) - (f*g)(x^{\prime}) \bigr| \le \|g\|_{p^{\prime}} \, \bigl\| (f - f_{-h})(x - \cdot) \bigr\|_p = \|g\|_{p^{\prime}} \, \|f - f_{-h}\|_p , \end{equation*}

where \(h = x^{\prime} - x\).

Step 3: continuity of translation in \(L^p(\mathbb{R})\). Because \(1 \le p < \infty\) and \(f \in L^p(\mathbb{R})\), Exercise 23(a) in Section 7A states that the function \(h \mapsto \|f - f_h\|_p\) is bounded and uniformly continuous on \(\mathbb{R}\). Since this function takes the value \(\|f - f_0\|_p = 0\) at \(h = 0\), its continuity at \(0\) gives

\begin{equation*} \lim_{h \to 0} \|f - f_{-h}\|_p = 0 . \end{equation*}

(For completeness: this limit can also be proved directly. By Exercise 18 in Section 7A there is, for each \(\varepsilon > 0\), a step function \(w \in L^p(\mathbb{R})\) with \(\|f - w\|_p < \varepsilon\). Because \(\|f - f_{-h}\|_p \le \|f - w\|_p + \|w - w_{-h}\|_p + \|w_{-h} - f_{-h}\|_p = 2\|f-w\|_p + \|w - w_{-h}\|_p\), it suffices to show \(\|w - w_{-h}\|_p \to 0\) for step functions \(w\), and by the triangle inequality it suffices to check this for \(w = \chi_{(a,b)}\) with \(a < b\), where \(\|w - w_{-h}\|_p^p \le 2|h|\) for all \(h\).)

Step 4: conclusion. Let \(\varepsilon > 0\). If \(g = 0\) then \(f * g = 0\) and there is nothing to prove, so assume \(\|g\|_{p^{\prime}} > 0\). By Step 3 there exists \(\delta > 0\) such that

\begin{equation*} \|f - f_{-h}\|_p < \frac{\varepsilon}{\|g\|_{p^{\prime}}} \quad \text{whenever } |h| < \delta . \end{equation*}

By Step 2, if \(x, x^{\prime} \in \mathbb{R}\) satisfy \(|x - x^{\prime}| < \delta\), then taking \(h = x^{\prime} - x\) we get

\begin{equation*} \bigl| (f*g)(x) - (f*g)(x^{\prime}) \bigr| \le \|g\|_{p^{\prime}} \, \|f - f_{-h}\|_p < \varepsilon . \end{equation*}

Because \(\delta\) depends only on \(\varepsilon\) and not on \(x\), the function \(f * g\) is uniformly continuous on \(\mathbb{R}\).

Problem (11C.11)

Suppose \(f \in L^\infty(\mathbb{R})\), \(x \in \mathbb{R}\), and \(f\) is continuous at \(x\). Prove that

\begin{equation*} \lim_{y \downarrow 0} (\mathcal{P}_y f)(x) = f(x) . \end{equation*}

Solution

Because \(f \in L^\infty(\mathbb{R})\) and \(P_y \in L^1(\mathbb{R})\) for each \(y > 0\), the function \(\mathcal{P}_y f = f * P_y\) is defined at every point of \(\mathbb{R}\) (see 11.70, with \(p = \infty\) and \(p^{\prime} = 1\)). By 11.65 we may write

\begin{equation*} (\mathcal{P}_y f)(x) = \int_{-\infty}^{\infty} P_y(t) \, f(x - t) \, dt ; \end{equation*}

the integrand is in \(L^1(\mathbb{R})\) as a function of \(t\) because \(|f(x-t)| \le \|f\|_\infty\) for almost every \(t\) and \(P_y \in L^1(\mathbb{R})\).

First note that \(|f(x)| \le \|f\|_\infty\). To see this, let \(\varepsilon > 0\). Continuity of \(f\) at \(x\) gives \(\delta > 0\) such that \(|f(z) - f(x)| < \varepsilon\) for all \(z\) with \(|z - x| < \delta\). The interval \((x - \delta, x + \delta)\) has positive Lebesgue measure, and \(|f| \le \|f\|_\infty\) almost everywhere, so there exists \(z\) in that interval with \(|f(z)| \le \|f\|_\infty\). Hence \(|f(x)| \le |f(z)| + \varepsilon \le \|f\|_\infty + \varepsilon\). Because \(\varepsilon > 0\) is arbitrary, \(|f(x)| \le \|f\|_\infty\). In particular

\begin{equation*} |f(x - t) - f(x)| \le 2\|f\|_\infty \quad \text{for almost every } t \in \mathbb{R} . \end{equation*}

Now let \(y > 0\). Because \(\int_{-\infty}^{\infty} P_y(t) \, dt = 1\) by 11.69(b), we have

\begin{equation*} (\mathcal{P}_y f)(x) - f(x) = \int_{-\infty}^{\infty} \bigl( f(x-t) - f(x) \bigr) P_y(t) \, dt , \end{equation*}

and hence

\begin{equation*} \bigl| (\mathcal{P}_y f)(x) - f(x) \bigr| \le \int_{-\infty}^{\infty} |f(x-t) - f(x)| \, P_y(t) \, dt , \end{equation*}

where we used \(P_y > 0\) (11.69(a)).

Let \(\varepsilon > 0\). Because \(f\) is continuous at \(x\), there exists \(\delta > 0\) such that

\begin{equation*} |f(x - t) - f(x)| < \varepsilon \quad \text{whenever } |t| < \delta . \end{equation*}

Split the last integral at \(|t| = \delta\). On \(\{t \in \mathbb{R} : |t| < \delta\}\) we get, using 11.69(a) and 11.69(b),

\begin{equation*} \int_{\{t : |t| < \delta\}} |f(x-t) - f(x)| \, P_y(t) \, dt \le \varepsilon \int_{-\infty}^{\infty} P_y(t) \, dt = \varepsilon . \end{equation*}

On \(\{t \in \mathbb{R} : |t| \ge \delta\}\) we get, using the bound \(|f(x-t) - f(x)| \le 2\|f\|_\infty\) valid for almost every \(t\),

\begin{equation*} \int_{\{t : |t| \ge \delta\}} |f(x-t) - f(x)| \, P_y(t) \, dt \le 2 \|f\|_\infty \int_{\{t : |t| \ge \delta\}} P_y(t) \, dt . \end{equation*}

Adding these two estimates gives

\begin{equation*} \bigl| (\mathcal{P}_y f)(x) - f(x) \bigr| \le \varepsilon + 2\|f\|_\infty \int_{\{t : |t| \ge \delta\}} P_y(t) \, dt \end{equation*}

for every \(y > 0\). By 11.69(c),

\begin{equation*} \lim_{y \downarrow 0} \int_{\{t \in \mathbb{R} \,:\, |t| \ge \delta\}} P_y(t) \, dt = 0 . \end{equation*}

Hence

\begin{equation*} \limsup_{y \downarrow 0} \bigl| (\mathcal{P}_y f)(x) - f(x) \bigr| \le \varepsilon . \end{equation*}

Because \(\varepsilon > 0\) was arbitrary, the limit superior above equals \(0\), which means

\begin{equation*} \lim_{y \downarrow 0} (\mathcal{P}_y f)(x) = f(x) , \end{equation*}

as desired.

Problem (11C.12)

Suppose \(p \in [1, \infty]\) and \(f \in L^p(\mathbb{R})\). Prove that \(\mathcal{P}_y(\mathcal{P}_{y^{\prime}} f) = \mathcal{P}_{y + y^{\prime}} f\) for all \(y, y^{\prime} > 0\).

Solution

Throughout, \(p^{\prime}\) denotes the exponent conjugate to \(p\).

Step 1: \(P_y \in L^q(\mathbb{R})\) for every \(y > 0\) and every \(q \in [1, \infty]\). Indeed, from 11.68 we have \(0 < P_y(x) = \frac{1}{\pi} \frac{y}{x^2 + y^2} \le \frac{1}{\pi y}\) for all \(x \in \mathbb{R}\), so \(\|P_y\|_\infty \le \frac{1}{\pi y}\), and \(\|P_y\|_1 = \int_{-\infty}^{\infty} P_y = 1\) by 11.69(a) and 11.69(b). Hence for \(q \in [1, \infty)\),

\begin{equation*} \|P_y\|_q^q = \int_{-\infty}^{\infty} P_y(x)^q \, dx \le \|P_y\|_\infty^{\,q-1} \int_{-\infty}^{\infty} P_y(x) \, dx = \|P_y\|_\infty^{\,q-1} < \infty . \end{equation*}

In particular \(P_y \in L^{p^{\prime}}(\mathbb{R})\), so by Hölder’s inequality (7.9) the number \((\mathcal{P}_y g)(x) = \int_{-\infty}^{\infty} g(t) P_y(x - t) \, dt\) is defined for every \(x \in \mathbb{R}\) whenever \(g \in L^p(\mathbb{R})\); moreover, by 11.64 and 11.65 (taking the \(L^1\) function to be \(P_y\)),

\begin{equation*} \|\mathcal{P}_y g\|_p = \|g * P_y\|_p \le \|P_y\|_1 \, \|g\|_p = \|g\|_p . \end{equation*}

Applying this with \(g = f\) shows \(\mathcal{P}_{y^{\prime}} f \in L^p(\mathbb{R})\), so \(\mathcal{P}_y(\mathcal{P}_{y^{\prime}} f)\) is defined at every point of \(\mathbb{R}\).

Step 2: the semigroup identity \(P_y * P_{y^{\prime}} = P_{y + y^{\prime}}\) for all \(y, y^{\prime} > 0\). Both \(P_y\) and \(P_{y^{\prime}}\) are in \(L^1(\mathbb{R})\), so \(P_y * P_{y^{\prime}} \in L^1(\mathbb{R})\) by 11.64 (with \(p = 1\)), and \(P_{y+y^{\prime}} \in L^1(\mathbb{R})\). By 11.66 and by 11.79 (which says \(\widehat{P_y}(t) = e^{-2\pi y |t|}\) for all \(t \in \mathbb{R}\)), for every \(t \in \mathbb{R}\) we have

\begin{equation*} \widehat{P_y * P_{y^{\prime}}}(t) = \widehat{P_y}(t) \, \widehat{P_{y^{\prime}}}(t) = e^{-2\pi y |t|} e^{-2\pi y^{\prime} |t|} = e^{-2\pi (y + y^{\prime}) |t|} = \widehat{P_{y + y^{\prime}}}(t) . \end{equation*}

Thus the function \(P_y * P_{y^{\prime}} - P_{y+y^{\prime}}\) belongs to \(L^1(\mathbb{R})\) and has Fourier transform identically \(0\); hence 11.80 implies that it is \(0\) as an element of \(L^1(\mathbb{R})\), meaning

\begin{equation*} (P_y * P_{y^{\prime}})(x) = P_{y + y^{\prime}}(x) \quad \text{for almost every } x \in \mathbb{R} . \end{equation*}

In fact this holds for every \(x \in \mathbb{R}\): the function \(P_{y+y^{\prime}}\) is continuous, and \(P_y * P_{y^{\prime}}\) is uniformly continuous on \(\mathbb{R}\) by Exercise 10 in this section (applied with the exponents \(1\) and \(\infty\), using \(P_y \in L^1(\mathbb{R})\) and \(P_{y^{\prime}} \in L^\infty(\mathbb{R})\) from Step 1), and two continuous functions on \(\mathbb{R}\) that agree almost everywhere agree everywhere.

Step 3: the associativity computation. Fix \(y, y^{\prime} > 0\) and \(x \in \mathbb{R}\), and write \(u = \mathcal{P}_{y^{\prime}} f\). Consider the function on \(\mathbb{R}^2\) defined by

\begin{equation*} F(s, t) = |f(t)| \, P_{y^{\prime}}(s - t) \, P_y(x - s) . \end{equation*}

This function is nonnegative and is measurable with respect to the product \(\sigma\)-algebra: the map \((s,t) \mapsto |f(t)|\) is product measurable because the inverse image of a Borel subset of \([0,\infty)\) under it is \(\mathbb{R} \times |f|^{-1}(\text{that set})\), and the maps \((s,t) \mapsto P_{y^{\prime}}(s-t)\) and \((s,t) \mapsto P_y(x-s)\) are continuous on \(\mathbb{R}^2\) and hence Borel measurable. Lebesgue measure on \(\mathbb{R}\) is \(\sigma\)-finite, so Tonelli’s Theorem (5.28) applies and gives

\begin{equation*} \int_{-\infty}^{\infty} \int_{-\infty}^{\infty} F(s,t) \, dt \, ds = \int_{-\infty}^{\infty} P_y(x - s) \, \bigl( |f| * P_{y^{\prime}} \bigr)(s) \, ds . \end{equation*}

By Step 1 applied to \(|f| \in L^p(\mathbb{R})\), we have \(|f| * P_{y^{\prime}} \in L^p(\mathbb{R})\) with \(\bigl\| |f| * P_{y^{\prime}} \bigr\|_p \le \|f\|_p\), and \(\|P_y(x - \cdot)\|_{p^{\prime}} = \|P_y\|_{p^{\prime}} < \infty\) (because \(P_y\) is an even function and Lebesgue measure is translation invariant by 2.7, the map \(s \mapsto P_y(x - s)\) has the same distribution of absolute values as \(P_y\)). Hence Hölder’s inequality (7.9) gives

\begin{equation*} \int_{-\infty}^{\infty} \int_{-\infty}^{\infty} F(s,t) \, dt \, ds \le \|P_y\|_{p^{\prime}} \, \|f\|_p < \infty . \end{equation*}

Thus Fubini’s Theorem (5.32), applied to the real and imaginary parts of the product measurable function \((s,t) \mapsto f(t) P_{y^{\prime}}(s-t) P_y(x-s)\), justifies the interchange of the order of integration below. Therefore

\begin{equation*} \begin{aligned} \bigl( \mathcal{P}_y (\mathcal{P}_{y^{\prime}} f) \bigr)(x) &= \int_{-\infty}^{\infty} u(s) \, P_y(x - s) \, ds \\ &= \int_{-\infty}^{\infty} \Bigl( \int_{-\infty}^{\infty} f(t) \, P_{y^{\prime}}(s - t) \, dt \Bigr) P_y(x - s) \, ds \\ &= \int_{-\infty}^{\infty} f(t) \Bigl( \int_{-\infty}^{\infty} P_{y^{\prime}}(s - t) \, P_y(x - s) \, ds \Bigr) dt . \end{aligned} \end{equation*}

For fixed \(t\), the substitution \(r = s - t\) (translation invariance of Lebesgue measure) shows that the inner integral equals

\begin{equation*} \int_{-\infty}^{\infty} P_{y^{\prime}}( r) \, P_y\bigl( (x - t) - r \bigr) \, dr = (P_{y^{\prime}} * P_y)(x - t) = P_{y + y^{\prime}}(x - t) , \end{equation*}

where the last equality is Step 2 (together with 11.65). Substituting this into the previous display gives

\begin{equation*} \bigl( \mathcal{P}_y (\mathcal{P}_{y^{\prime}} f) \bigr)(x) = \int_{-\infty}^{\infty} f(t) \, P_{y + y^{\prime}}(x - t) \, dt = \bigl( \mathcal{P}_{y + y^{\prime}} f \bigr)(x) . \end{equation*}

Because \(x \in \mathbb{R}\) was arbitrary, \(\mathcal{P}_y(\mathcal{P}_{y^{\prime}} f) = \mathcal{P}_{y+y^{\prime}} f\), as desired.

Problem (11C.13)

Suppose \(p \in [1, \infty]\) and \(f \in L^p(\mathbb{R})\). Prove that if \(0 < y < y^{\prime}\), then

\begin{equation*} \|\mathcal{P}_y f\|_p \ge \|\mathcal{P}_{y^{\prime}} f\|_p . \end{equation*}

Solution

First we record the basic norm estimate. If \(g \in L^p(\mathbb{R})\) and \(s > 0\), then \(P_s \in L^1(\mathbb{R})\) with \(\|P_s\|_1 = \int_{-\infty}^{\infty} P_s(x)\,dx = 1\) by 11.69(a) and 11.69(b). Hence 11.64 (applied with the \(L^1\) function \(P_s\) and the \(L^p\) function \(g\)) together with 11.65 gives \(\mathcal{P}_s g = g * P_s = P_s * g \in L^p(\mathbb{R})\) and

\begin{equation*} \|\mathcal{P}_s g\|_p = \|P_s * g\|_p \le \|P_s\|_1 \, \|g\|_p = \|g\|_p . \end{equation*}

In words: each operator \(\mathcal{P}_s\) maps \(L^p(\mathbb{R})\) into \(L^p(\mathbb{R})\) and does not increase the \(L^p\)-norm.

Now suppose \(0 < y < y^{\prime}\), and set \(s = y^{\prime} - y\), so \(s > 0\) and \(y^{\prime} = s + y\). By Exercise 12 in this section (applied with \(s\) in place of \(y\) and \(y\) in place of \(y^{\prime}\)),

\begin{equation*} \mathcal{P}_{y^{\prime}} f = \mathcal{P}_{s + y} f = \mathcal{P}_s (\mathcal{P}_y f) . \end{equation*}

By the previous paragraph, \(\mathcal{P}_y f \in L^p(\mathbb{R})\), so applying the norm estimate with \(g = \mathcal{P}_y f\) gives

\begin{equation*} \|\mathcal{P}_{y^{\prime}} f\|_p = \|\mathcal{P}_s (\mathcal{P}_y f)\|_p \le \|\mathcal{P}_y f\|_p . \end{equation*}

That is, \(\|\mathcal{P}_y f\|_p \ge \|\mathcal{P}_{y^{\prime}} f\|_p\), as desired.

Problem (11C.14)

Suppose \(f \in L^1(\mathbb{R})\).

(a) Prove that \(\widehat{\overline{f}}(t) = \overline{\hat{f}(-t)}\) for all \(t \in \mathbb{R}\).

(b) Prove that \(f(x) \in \mathbb{R}\) for almost every \(x \in \mathbb{R}\) if and only if \(\hat{f}(t) = \overline{\hat{f}(-t)}\) for all \(t \in \mathbb{R}\).

Solution

Here \(\overline{f}\) denotes the function \(x \mapsto \overline{f(x)}\). Because \(|\overline{f}| = |f|\), we have \(\overline{f} \in L^1(\mathbb{R})\), so \(\widehat{\overline{f}}\) is defined.

(a) Fix \(t \in \mathbb{R}\). Complex conjugation commutes with the integral (apply the definition of the integral of a complex-valued function to real and imaginary parts, or note that \(\overline{\int h} = \int \overline{h}\) for \(h \in L^1\)). Because \(\overline{e^{-2\pi i t x}} = e^{2\pi i t x} = e^{-2\pi i (-t) x}\) for all \(x \in \mathbb{R}\), we obtain

\begin{equation*} \begin{aligned} \widehat{\overline{f}}(t) &= \int_{-\infty}^{\infty} \overline{f(x)} \, e^{-2\pi i t x} \, dx \\ &= \int_{-\infty}^{\infty} \overline{ f(x) \, e^{2\pi i t x} } \, dx \\ &= \overline{ \int_{-\infty}^{\infty} f(x) \, e^{2\pi i t x} \, dx } \\ &= \overline{ \int_{-\infty}^{\infty} f(x) \, e^{-2\pi i (-t) x} \, dx } \\ &= \overline{ \hat{f}(-t) } . \end{aligned} \end{equation*}

(b) First suppose \(f(x) \in \mathbb{R}\) for almost every \(x \in \mathbb{R}\). Then \(\overline{f(x)} = f(x)\) for almost every \(x \in \mathbb{R}\), so \(\overline{f}\) and \(f\) are equal as elements of \(L^1(\mathbb{R})\). Because functions that agree almost everywhere have the same integral against \(e^{-2\pi i t x}\), this gives \(\widehat{\overline{f}} = \hat{f}\). Combining this with part (a), we conclude that

\begin{equation*} \hat{f}(t) = \widehat{\overline{f}}(t) = \overline{\hat{f}(-t)} \end{equation*}

for all \(t \in \mathbb{R}\).

Conversely, suppose \(\hat{f}(t) = \overline{\hat{f}(-t)}\) for all \(t \in \mathbb{R}\). By part (a), \(\overline{\hat{f}(-t)} = \widehat{\overline{f}}(t)\), so

\begin{equation*} \hat{f}(t) = \widehat{\overline{f}}(t) \end{equation*}

for all \(t \in \mathbb{R}\). The function \(f - \overline{f}\) is in \(L^1(\mathbb{R})\), and the Fourier transform is linear, so

\begin{equation*} \widehat{(f - \overline{f})}(t) = \hat{f}(t) - \widehat{\overline{f}}(t) = 0 \end{equation*}

for all \(t \in \mathbb{R}\). Hence 11.80 implies that \(f - \overline{f} = 0\) as an element of \(L^1(\mathbb{R})\); in other words,

\begin{equation*} f(x) = \overline{f(x)} \quad \text{for almost every } x \in \mathbb{R} . \end{equation*}

A complex number equals its own complex conjugate if and only if it is real, so \(f(x) \in \mathbb{R}\) for almost every \(x \in \mathbb{R}\), completing the proof.

Problem (11C.15)

Define \(f \in L^1(\mathbb{R})\) by \(f(x) = e^{-x^4} \chi_{[0,\infty)}(x)\). Show that \(\widehat{f} \notin L^1(\mathbb{R})\).

Solution

First note that \(f\) really is in \(L^1(\mathbb{R})\): the function \(f\) is Borel measurable, \(0 \le f \le 1\), and

\begin{equation*} \|f\|_1 = \int_0^{\infty} e^{-x^4}\,dx \le \int_0^{1} 1\,dx + \int_1^{\infty} e^{-x}\,dx < \infty , \end{equation*}

where we used \(x^4 \ge x\) for \(x \ge 1\).

Now suppose, seeking a contradiction, that \(\widehat{f} \in L^1(\mathbb{R})\). Then \(f\) and \(\widehat{f}\) are both in \(L^1(\mathbb{R})\), so the Fourier Inversion Formula (11.76) applies. Define \(h : \mathbb{R} \to \mathbb{C}\) by

\begin{equation*} h(x) = \int_{-\infty}^{\infty} \widehat{f}(t)\, e^{2\pi i x t}\,dt . \end{equation*}

The integral above makes sense for every \(x \in \mathbb{R}\) because \(|\widehat{f}(t) e^{2\pi i x t}| = |\widehat{f}(t)|\) and \(\widehat{f} \in L^1(\mathbb{R})\). Furthermore \(h(x) = (\widehat{f}\,)\widehat{\phantom{f}}(-x)\), so \(h\) is continuous on \(\mathbb{R}\) by 11.49 applied to the function \(\widehat{f} \in L^1(\mathbb{R})\). (Alternatively: if \(x_n \to x\), then \(\widehat{f}(t) e^{2\pi i x_n t} \to \widehat{f}(t) e^{2\pi i x t}\) for every \(t\), with all these functions dominated in absolute value by \(|\widehat{f}| \in L^1(\mathbb{R})\); the Dominated Convergence Theorem (3.31) then gives \(h(x_n) \to h(x)\).)

The Fourier Inversion Formula (11.76) states that \(f(x) = h(x)\) for almost every \(x \in \mathbb{R}\). Let

\begin{equation*} E = \{x \in \mathbb{R} : f(x) = h(x)\} . \end{equation*}

Thus \(|\mathbb{R} \setminus E| = 0\).

Fix \(n \in \mathbb{Z}^+\). The interval \((0, \tfrac1n)\) has Lebesgue measure \(\tfrac1n > 0\), so it is not contained in the null set \(\mathbb{R} \setminus E\); hence there exists \(a_n \in (0, \tfrac1n) \cap E\). Similarly there exists \(b_n \in (-\tfrac1n, 0) \cap E\).

Because \(a_n \in (0, \tfrac1n)\) we have \(a_n \to 0\), and because \(a_n > 0\) we have

\begin{equation*} h(a_n) = f(a_n) = e^{-a_n^{\,4}} \longrightarrow 1 \end{equation*}

as \(n \to \infty\). The continuity of \(h\) at \(0\) therefore forces \(h(0) = 1\).

Because \(b_n \in (-\tfrac1n, 0)\) we have \(b_n \to 0\), and because \(b_n < 0\) we have \(f(b_n) = 0\), so

\begin{equation*} h(b_n) = f(b_n) = 0 \longrightarrow 0 . \end{equation*}

The continuity of \(h\) at \(0\) therefore forces \(h(0) = 0\).

We have derived \(1 = h(0) = 0\), a contradiction. Hence the assumption that \(\widehat{f} \in L^1(\mathbb{R})\) is untenable, and thus \(\widehat{f} \notin L^1(\mathbb{R})\).

In words: if \(\widehat{f}\) were in \(L^1(\mathbb{R})\), then \(f\) would agree almost everywhere with a continuous function, which is impossible because \(f\) has a genuine jump of size \(1\) at the origin (it is continuous and equal to \(0\) on \((-\infty, 0)\) and continuous with limit \(1\) from the right at \(0\)).

Problem (11C.16)

Suppose \(f \in L^1(\mathbb{R})\) and \(\widehat{f} \in L^1(\mathbb{R})\). Prove that \(f \in L^2(\mathbb{R})\) and \(\widehat{f} \in L^2(\mathbb{R})\).

Solution

The key point is that the hypotheses force \(f\) to be (almost everywhere equal to) a bounded function; a bounded function in \(L^1(\mathbb{R})\) is automatically in \(L^2(\mathbb{R})\).

Step 1: \(f \in L^{\infty}(\mathbb{R})\).

Because \(f \in L^1(\mathbb{R})\) and \(\widehat{f} \in L^1(\mathbb{R})\), the Fourier Inversion Formula (11.76) gives

\begin{equation*} f(x) = \bigl(\widehat{f}\,\bigr)\widehat{\phantom{f}}(-x) \end{equation*}

for almost every \(x \in \mathbb{R}\). Now \(\widehat{f} \in L^1(\mathbb{R})\), so 11.49 applied to \(\widehat{f}\) shows that \((\widehat{f}\,)\widehat{\phantom{f}}\) is a bounded function on \(\mathbb{R}\); indeed, directly from the definition of the Fourier transform,

\begin{equation*} \bigl|\bigl(\widehat{f}\,\bigr)\widehat{\phantom{f}}(s)\bigr| = \left| \int_{-\infty}^{\infty} \widehat{f}(t)\, e^{-2\pi i s t}\,dt \right| \le \int_{-\infty}^{\infty} |\widehat{f}(t)|\,dt = \|\widehat{f}\,\|_1 \end{equation*}

for every \(s \in \mathbb{R}\). Combining the two displays above, we get

\begin{equation*} |f(x)| \le \|\widehat{f}\,\|_1 \quad \text{for almost every } x \in \mathbb{R}, \end{equation*}

which means \(f \in L^{\infty}(\mathbb{R})\) with \(\|f\|_{\infty} \le \|\widehat{f}\,\|_1 < \infty\).

Step 2: \(f \in L^2(\mathbb{R})\).

Using Step 1 and the hypothesis \(f \in L^1(\mathbb{R})\), we have

\begin{equation*} \|f\|_2^2 = \int_{-\infty}^{\infty} |f(x)|^2\,dx = \int_{-\infty}^{\infty} |f(x)|\,|f(x)|\,dx \le \|f\|_{\infty} \int_{-\infty}^{\infty} |f(x)|\,dx = \|f\|_{\infty}\,\|f\|_1 < \infty . \end{equation*}

Hence \(f \in L^2(\mathbb{R})\). (This is the standard interpolation fact: \(L^1(\mathbb{R}) \cap L^{\infty}(\mathbb{R}) \subseteq L^p(\mathbb{R})\) for all \(p \in [1,\infty]\).)

Step 3: \(\widehat{f} \in L^2(\mathbb{R})\).

By Step 2 together with the hypothesis, \(f \in L^1(\mathbb{R}) \cap L^2(\mathbb{R})\). Thus Plancherel’s Theorem (11.82) applies and gives

\begin{equation*} \|\widehat{f}\,\|_2 = \|f\|_2 < \infty . \end{equation*}

Hence \(\widehat{f} \in L^2(\mathbb{R})\), completing the proof.

Remark: symmetry gives a second route to Step 3. Set \(g = \widehat{f}\). Then \(g \in L^1(\mathbb{R})\) and, by the Fourier Inversion Formula as used in Step 1, \(\widehat{g} = (\widehat{f}\,)\widehat{\phantom{f}}\) agrees almost everywhere with \(x \mapsto f(-x)\), which is in \(L^1(\mathbb{R})\). So the pair \((g, \widehat{g})\) satisfies the same hypotheses as \((f, \widehat{f}\,)\), and Steps 1 and 2 applied to \(g\) give \(g = \widehat{f} \in L^2(\mathbb{R})\) without invoking Plancherel’s Theorem.

Problem (11C.17)

Prove there exists a continuous function \(g : \mathbb{R} \to \mathbb{R}\) such that \(\lim\limits_{t \to \pm\infty} g(t) = 0\) and \(g \notin \{\widehat{f} : f \in L^1(\mathbb{R})\}\).

Solution

By 11.49, every function of the form \(\widehat{f}\) with \(f \in L^1(\mathbb{R})\) is continuous and has limit \(0\) at \(\pm\infty\). This exercise asks us to show that the converse fails, i.e. that the Fourier transform does not map \(L^1(\mathbb{R})\) onto the space of continuous functions vanishing at \(\pm\infty\).

The function. Define \(g : \mathbb{R} \to \mathbb{R}\) by

\begin{equation*} g(t) = \begin{cases} \dfrac{1}{\log t} & \text{if } t \ge 2, \\ \dfrac{t}{2 \log 2} & \text{if } 0 \le t < 2, \\ -g(-t) & \text{if } t < 0 . \end{cases} \end{equation*}

Then \(g\) is real valued; \(g\) is continuous on \([2, \infty)\) and on \([0,2]\), and the two formulas agree at \(t = 2\) (both give \(\frac{1}{\log 2}\)), so \(g\) is continuous on \([0,\infty)\); because \(g(0) = 0\) and \(g\) is odd, \(g\) is continuous on all of \(\mathbb{R}\). Also \(g(t) = \frac{1}{\log t} \to 0\) as \(t \to \infty\), and hence \(g(t) \to 0\) as \(t \to -\infty\) by oddness. So \(g\) is continuous, real valued, odd, and vanishes at \(\pm \infty\).

The crucial property of \(g\) is that

\begin{equation*} \int_2^{\infty} \frac{g(t)}{t}\,dt = \int_2^{\infty} \frac{dt}{t \log t} = \lim_{T \to \infty}\bigl( \log \log T - \log \log 2 \bigr) = \infty , \end{equation*}

while \(g(t)/t\) is bounded on \([1,2]\). Hence

\begin{equation*} \lim_{T \to \infty} \int_1^T \frac{g(t)}{t}\,dt = \infty . \end{equation*}

Call this last statement the divergence property of \(g\).

A lemma on the sine integral. Define \(S : [0,\infty) \to \mathbb{R}\) by \(S(A) = \int_0^A \frac{\sin u}{u}\,du\), where the integrand is given the value \(1\) at \(u = 0\) so that it is continuous on \([0,\infty)\). We claim that

\begin{equation*} M := \sup_{A \ge 0} |S(A)| < \infty \qquad \text{and} \qquad L := \lim_{A \to \infty} S(A) \text{ exists in } \mathbb{R}. \end{equation*}

Indeed, \(|S(A)| \le 1\) for \(A \in [0,1]\) because the integrand has absolute value at most \(1\) there. For \(A > 1\), integration by parts gives

\begin{equation*} \int_1^A \frac{\sin u}{u}\,du = \Bigl[ \frac{-\cos u}{u} \Bigr]_{u=1}^{u=A} - \int_1^A \frac{\cos u}{u^2}\,du = \cos 1 - \frac{\cos A}{A} - \int_1^A \frac{\cos u}{u^2}\,du , \end{equation*}

whose absolute value is at most \(1 + 1 + \int_1^{\infty} u^{-2}\,du = 3\). Hence \(|S(A)| \le 4\) for all \(A \ge 0\). Moreover \(\frac{\cos A}{A} \to 0\) as \(A \to \infty\) and \(\int_1^{\infty} \frac{\cos u}{u^2}\,du\) converges absolutely, so the displayed equation shows that \(\lim_{A \to \infty} S(A)\) exists in \(\mathbb{R}\). This proves the claim.

Now suppose, seeking a contradiction, that there exists \(f \in L^1(\mathbb{R})\) with \(\widehat{f} = g\).

Step 1: reduction to a real odd function. Define \(f_1 : \mathbb{R} \to \mathbb{C}\) by \(f_1(x) = \frac{f(x) - f(-x)}{2}\). Then \(f_1 \in L^1(\mathbb{R})\) and \(f_1\) is odd. Because \(\int_{-\infty}^{\infty} f(-x) e^{-2\pi i t x}\,dx = \int_{-\infty}^{\infty} f(y) e^{2\pi i t y}\,dy = \widehat{f}(-t)\) (substitute \(y = -x\), using the reflection invariance of Lebesgue measure), we get

\begin{equation*} \widehat{f_1}(t) = \frac{\widehat{f}(t) - \widehat{f}(-t)}{2} = \frac{g(t) - g(-t)}{2} = g(t) \end{equation*}

for all \(t \in \mathbb{R}\), the last equality because \(g\) is odd.

Write \(f_1 = u + i v\), where \(u = \operatorname{Re} f_1\) and \(v = \operatorname{Im} f_1\). Then \(u, v \in L^1(\mathbb{R})\) are real valued and odd.

If \(w \in L^1(\mathbb{R})\) is real valued and odd and \(t \in \mathbb{R}\), then

\begin{equation*} \widehat{w}(t) = \int_{-\infty}^{\infty} w(x) \cos(2\pi t x)\,dx - i \int_{-\infty}^{\infty} w(x) \sin(2 \pi t x)\,dx = -2 i \int_0^{\infty} w(x) \sin(2\pi t x)\,dx , \end{equation*}

because \(x \mapsto w(x)\cos(2\pi t x)\) is an odd integrable function (hence has integral \(0\) over \(\mathbb{R}\)) and \(x \mapsto w(x) \sin(2\pi t x)\) is an even integrable function. In particular \(\widehat{w}(t)\) is purely imaginary.

Applying this to \(u\) and to \(v\) and writing \(A(t) = \int_0^{\infty} u(x) \sin(2 \pi t x)\,dx\) and \(B(t) = \int_0^{\infty} v(x) \sin(2\pi t x)\,dx\) (both real numbers), we obtain

\begin{equation*} g(t) = \widehat{f_1}(t) = \widehat{u}(t) + i \widehat{v}(t) = -2 i A(t) + 2 B(t) . \end{equation*}

Because \(g\) is real valued, taking real parts gives

\begin{equation*} g(t) = 2 \int_0^{\infty} v(x) \sin(2\pi t x)\,dx \qquad \text{for all } t \in \mathbb{R}, \end{equation*}

where \(v \in L^1(\mathbb{R})\) is real valued.

Step 2: the integral must converge. Fix \(T > 1\). Then

\begin{equation*} \int_1^T \frac{g(t)}{t}\,dt = 2 \int_1^T \frac{1}{t} \int_0^{\infty} v(x) \sin(2 \pi t x)\,dx\,dt . \end{equation*}

The function \((t, x) \mapsto \frac{v(x) \sin(2 \pi t x)}{t}\) is measurable on \([1,T] \times (0,\infty)\) (it is the product of a continuous function of \((t,x)\) with a measurable function of \(x\)), and both factor spaces carry Lebesgue measure, which is \(\sigma\)-finite. Because \(|\sin| \le 1\) and \(\int_0^{\infty} |v| \le \|v\|_1\), Tonelli’s Theorem (5.28) gives

\begin{equation*} \int_1^T \int_0^{\infty} \frac{|v(x)| \, |\sin(2 \pi t x)|}{t}\,dx\,dt \le \int_1^T \frac{\|v\|_1}{t}\,dt = \|v\|_1 \log T < \infty . \end{equation*}

Thus Fubini’s Theorem (5.32) applies to the iterated integral above, and hence

\begin{equation*} \int_1^T \frac{g(t)}{t}\,dt = 2 \int_0^{\infty} v(x) \left( \int_1^T \frac{\sin(2\pi t x)}{t}\,dt \right) dx . \end{equation*}

For \(x > 0\) the substitution \(u = 2 \pi t x\) (so that \(\frac{dt}{t} = \frac{du}{u}\)) gives

\begin{equation*} \int_1^T \frac{\sin(2 \pi t x)}{t}\,dt = \int_{2 \pi x}^{2 \pi T x} \frac{\sin u}{u}\,du = S(2 \pi T x) - S(2 \pi x) =: \Phi_T(x) . \end{equation*}

By the lemma, \(|\Phi_T(x)| \le 2M\) for all \(x > 0\) and all \(T > 1\), and for each fixed \(x > 0\) we have \(\Phi_T(x) \to L - S(2\pi x)\) as \(T \to \infty\).

Take any sequence \(T_1, T_2, \ldots\) in \((1, \infty)\) with \(T_n \to \infty\). The functions \(x \mapsto v(x) \Phi_{T_n}(x)\) converge pointwise on \((0,\infty)\) to \(x \mapsto v(x)\bigl(L - S(2\pi x)\bigr)\) and are dominated in absolute value by \(2M |v| \in L^1(\mathbb{R})\). Hence the Dominated Convergence Theorem (3.31) gives

\begin{equation*} \lim_{n \to \infty} \int_1^{T_n} \frac{g(t)}{t}\,dt = 2 \int_0^{\infty} v(x)\bigl(L - S(2 \pi x)\bigr)\,dx , \end{equation*}

and the right side is a real number (its absolute value is at most \(4M \|v\|_1\)).

Step 3: contradiction. Step 2 shows that \(\int_1^{T_n} \frac{g(t)}{t}\,dt\) converges to a finite limit for every sequence \(T_n \to \infty\), contradicting the divergence property of \(g\), which says this quantity tends to \(\infty\).

Hence no \(f \in L^1(\mathbb{R})\) satisfies \(\widehat{f} = g\). In other words, \(g\) is a continuous real-valued function on \(\mathbb{R}\) with \(\lim_{t \to \pm\infty} g(t) = 0\) such that \(g \notin \{\widehat{f} : f \in L^1(\mathbb{R})\}\), as desired.

Remark: the proof shows more, namely that if \(g\) is odd, real valued, and lies in the range of the Fourier transform on \(L^1(\mathbb{R})\), then \(\lim_{T \to \infty} \int_1^T \frac{g(t)}{t}\,dt\) exists as a real number. No such constraint is imposed on a continuous function vanishing at \(\pm\infty\), and hence \(\{\widehat{f} : f \in L^1(\mathbb{R})\}\) is a proper subset of the space of continuous functions on \(\mathbb{R}\) with limit \(0\) at \(\pm\infty\).

Problem (11C.18)

Prove that if \(f \in L^1(\mathbb{R})\), then \(\|\widehat{f}\,\|_2 = \|f\|_2\).

[This exercise slightly improves Plancherel’s Theorem (11.82) because here we have the weaker hypothesis that \(f \in L^1(\mathbb{R})\) instead of \(f \in L^1(\mathbb{R}) \cap L^2(\mathbb{R})\). Because of Plancherel’s Theorem, here you need only prove that if \(f \in L^1(\mathbb{R})\) and \(\|f\|_2 = \infty\), then \(\|\widehat{f}\,\|_2 = \infty\).]

Solution

Here \(\|h\|_2\) denotes \(\bigl(\int_{-\infty}^{\infty} |h|^2\bigr)^{1/2} \in [0,\infty]\) for an arbitrary measurable \(h\), so the asserted equality includes the case where both sides equal \(\infty\).

Suppose \(f \in L^1(\mathbb{R})\).

Case 1: \(\|f\|_2 < \infty\). Then \(f \in L^1(\mathbb{R}) \cap L^2(\mathbb{R})\), so Plancherel’s Theorem (11.82) gives \(\|\widehat{f}\,\|_2 = \|f\|_2\), as desired.

Case 2: \(\|f\|_2 = \infty\). We must show \(\|\widehat{f}\,\|_2 = \infty\).

For \(y > 0\) consider \(f * P_y\), where \(P_y\) is the Poisson kernel on the upper half-plane (11.68). We record three facts.

(i) \(f * P_y \in L^1(\mathbb{R}) \cap L^2(\mathbb{R})\). First, \(P_y \in L^1(\mathbb{R})\) with \(\|P_y\|_1 = 1\) by 11.69(a) and 11.69(b). Also \(P_y \in L^2(\mathbb{R})\), because \(0 < P_y(x) \le \frac{1}{\pi y}\) for all \(x\) and \(P_y(x) \le \frac{y}{\pi x^2}\) for \(x \ne 0\), so \(\int_{-\infty}^{\infty} P_y^2 \le \int_{-1}^{1} \frac{1}{\pi^2 y^2} + \int_{|x| \ge 1} \frac{y^2}{\pi^2 x^4} < \infty\). Now apply 11.64 twice, each time with the \(L^1\) function \(f\) in the first slot: with \(p = 1\) and \(g = P_y \in L^1(\mathbb{R})\) it gives \(\|f * P_y\|_1 \le \|f\|_1 \|P_y\|_1 < \infty\), and with \(p = 2\) and \(g = P_y \in L^2(\mathbb{R})\) it gives \(\|f * P_y\|_2 \le \|f\|_1 \|P_y\|_2 < \infty\).

(ii) \((f * P_y)\widehat{\phantom{f}}(t) = \widehat{f}(t)\, e^{-2\pi y |t|}\) for all \(t \in \mathbb{R}\). This is 11.85: it follows from 11.66 (the Fourier transform of a convolution is the product of the Fourier transforms) together with 11.79 (which says \((P_y)\widehat{\phantom{f}}(t) = e^{-2\pi y |t|}\)).

(iii) By (i), Plancherel’s Theorem (11.82) applies to \(f * P_y\), so by (ii)

\begin{equation*} \|f * P_y\|_2 = \|(f * P_y)\widehat{\phantom{f}}\,\|_2 = \left( \int_{-\infty}^{\infty} |\widehat{f}(t)|^2 e^{-4 \pi y |t|}\,dt \right)^{1/2} \le \|\widehat{f}\,\|_2 , \end{equation*}

where the last inequality holds because \(e^{-4\pi y|t|} \le 1\).

Now 11.74 (with \(p = 1\)) states that \(\lim_{y \downarrow 0} \|f - f * P_y\|_1 = 0\). Choose a sequence \(y_1, y_2, \ldots\) of positive numbers with \(y_n \to 0\); then \(\|f - f*P_{y_n}\|_1 \to 0\), so by 7.23 there is a subsequence, which after relabeling we again call \(y_1, y_2, \ldots\), such that

\begin{equation*} (f * P_{y_n})(x) \to f(x) \quad \text{for almost every } x \in \mathbb{R}. \end{equation*}

Hence \(|f(x)|^2 = \lim_{n \to \infty} |(f * P_{y_n})(x)|^2\) for almost every \(x\), and Fatou’s Lemma (Exercise 17 in Section 3A, applied to the nonnegative measurable functions \(|f * P_{y_n}|^2\)) gives

\begin{equation*} \|f\|_2^2 = \int_{-\infty}^{\infty} \liminf_{n \to \infty} |(f * P_{y_n})(x)|^2\,dx \le \liminf_{n \to \infty} \int_{-\infty}^{\infty} |(f * P_{y_n})(x)|^2\,dx = \liminf_{n \to \infty} \|f * P_{y_n}\|_2^2 . \end{equation*}

By (iii), \(\|f * P_{y_n}\|_2^2 \le \|\widehat{f}\,\|_2^2\) for every \(n\). Combining the two previous sentences yields

\begin{equation*} \|f\|_2^2 \le \|\widehat{f}\,\|_2^2 . \end{equation*}

Because \(\|f\|_2 = \infty\) in this case, we conclude that \(\|\widehat{f}\,\|_2 = \infty = \|f\|_2\), as desired.

The two cases together prove that \(\|\widehat{f}\,\|_2 = \|f\|_2\) for every \(f \in L^1(\mathbb{R})\).

Remark: nothing in the derivation of \(\|f\|_2 \le \|\widehat{f}\,\|_2\) in Case 2 used the assumption \(\|f\|_2 = \infty\), so that inequality holds for every \(f \in L^1(\mathbb{R})\). The reverse inequality holds for every \(f \in L^1(\mathbb{R})\) too, and again without invoking Case 1. Indeed, the inequality is trivial when \(\|f\|_2 = \infty\); and if \(\|f\|_2 < \infty\), then 11.65 and 11.64 (applied with \(P_y\) in the \(L^1\) slot and \(f \in L^2(\mathbb{R})\)) give \(\|f * P_y\|_2 = \|P_y * f\|_2 \le \|P_y\|_1 \|f\|_2 = \|f\|_2\), while for any decreasing sequence \(y_1 > y_2 > \cdots\) with \(y_n \to 0\) the functions \(|\widehat{f}\,|^2 e^{-4\pi y_n |t|}\) increase pointwise to \(|\widehat{f}\,|^2\), so the Monotone Convergence Theorem (3.11) shows that the middle quantity in (iii) increases to \(\|\widehat{f}\,\|_2\). Thus the equality can be obtained without splitting into cases.

Problem (11C.19)

Suppose \(y > 0\). Define an operator \(T\) on \(L^2(\mathbb{R})\) by \(Tf = f * P_y\).

(a) Show that \(T\) is a self-adjoint operator on \(L^2(\mathbb{R})\).

(b) Show that \(\operatorname{sp}(T) = [0,1]\).

[Because the spectrum of each compact operator is a countable set (by 10.93), part (b) above implies that \(T\) is not a compact operator. This conclusion differs from the situation on the unit circle—see Exercise 9 in Section 11B.]

Solution

Throughout, \(P_y(x) = \frac{1}{\pi} \frac{y}{x^2 + y^2}\) is the Poisson kernel on the upper half-plane (11.68). Two elementary properties of \(P_y\) are used repeatedly: \(P_y\) is real valued and even (obvious from the formula, since \(x\) appears only through \(x^2\)), and \(\|P_y\|_1 = 1\) by 11.69(a) and 11.69(b).

\(T\) is a bounded operator. If \(f \in L^2(\mathbb{R})\), then by 11.65 and 11.64 (applied with the \(L^1\) function \(P_y\) and the \(L^2\) function \(f\)),

\begin{equation*} \|Tf\|_2 = \|f * P_y\|_2 = \|P_y * f\|_2 \le \|P_y\|_1 \|f\|_2 = \|f\|_2 . \end{equation*}

In particular \(Tf \in L^2(\mathbb{R})\), and \(T\) is a bounded linear operator on \(L^2(\mathbb{R})\) with \(\|T\| \le 1\). (Linearity of \(T\) is clear from the definition of convolution.)

(a) Let \(f, g \in L^2(\mathbb{R})\). Then, using the inner product \(\langle h_1, h_2 \rangle = \int_{-\infty}^{\infty} h_1 \overline{h_2}\) on \(L^2(\mathbb{R})\),

\begin{equation*} \langle Tf, g \rangle = \int_{-\infty}^{\infty} \left( \int_{-\infty}^{\infty} f(t) P_y(x - t)\,dt \right) \overline{g(x)}\,dx . \end{equation*}

The function \((t, x) \mapsto f(t) P_y(x - t) \overline{g(x)}\) is measurable on \(\mathbb{R}^2\): the maps \((t,x) \mapsto f(t)\) and \((t,x) \mapsto \overline{g(x)}\) are measurable because they are compositions of measurable functions with the (continuous, hence Borel measurable) coordinate projections, and \((t,x) \mapsto P_y(x - t)\) is continuous. Lebesgue measure on \(\mathbb{R}\) is \(\sigma\)-finite, so Tonelli’s Theorem (5.28) applies to the absolute value of this function and gives, using 11.65, 11.64, and Hölder’s inequality (7.9),

\begin{equation*} \int_{-\infty}^{\infty} \int_{-\infty}^{\infty} |f(t)| P_y(x-t) |g(x)|\,dt\,dx = \int_{-\infty}^{\infty} \bigl( |f| * P_y \bigr)(x)\, |g(x)|\,dx \le \bigl\| |f| * P_y \bigr\|_2 \|g\|_2 \le \|f\|_2 \|g\|_2 < \infty . \end{equation*}

Hence the double integral converges absolutely and Fubini’s Theorem (5.32) permits interchanging the order of integration:

\begin{equation*} \langle Tf, g \rangle = \int_{-\infty}^{\infty} f(t) \left( \int_{-\infty}^{\infty} P_y(x - t) \overline{g(x)}\,dx \right) dt . \end{equation*}

Because \(P_y\) is real valued and even, \(P_y(x - t) = P_y(t - x)\), so the inner integral equals

\begin{equation*} \int_{-\infty}^{\infty} P_y(t - x) \overline{g(x)}\,dx = \overline{ \int_{-\infty}^{\infty} g(x) P_y(t-x)\,dx } = \overline{(g * P_y)(t)} = \overline{(Tg)(t)} . \end{equation*}

Therefore

\begin{equation*} \langle Tf, g \rangle = \int_{-\infty}^{\infty} f(t) \overline{(Tg)(t)}\,dt = \langle f, Tg \rangle \end{equation*}

for all \(f, g \in L^2(\mathbb{R})\). Thus \(T^* = T\), so \(T\) is self-adjoint.

(b) Let \(m : \mathbb{R} \to (0, 1]\) be defined by \(m(t) = e^{-2\pi y |t|}\), and let \(M_m\) denote the multiplication operator on \(L^2(\mathbb{R})\) defined by \(M_m h = m h\). Because \(|m| \le 1\), the operator \(M_m\) is bounded with \(\|M_m\| \le 1\).

Step 1: \(T\) is unitarily equivalent to \(M_m\).

Let \(f \in L^2(\mathbb{R})\). Since \(P_y \in L^1(\mathbb{R})\), Exercise 20 in this section (together with 11.65) gives

\begin{equation*} \mathcal{F}(Tf) = \mathcal{F}(P_y * f) = \widehat{P_y} \cdot \mathcal{F} f = m \cdot \mathcal{F} f = M_m (\mathcal{F} f), \end{equation*}

where the third equality uses 11.79, which says \(\widehat{P_y}(t) = e^{-2 \pi y |t|}\) for all \(t \in \mathbb{R}\). Thus \(\mathcal{F} T = M_m \mathcal{F}\). By 11.87(a) the Fourier transform \(\mathcal{F}\) is a unitary operator on \(L^2(\mathbb{R})\), hence invertible, and therefore

\begin{equation*} T = \mathcal{F}^{-1} M_m \mathcal{F} . \end{equation*}

(As a check on part (a): \(T^* = \mathcal{F}^{-1} M_{\overline{m}} \mathcal{F} = \mathcal{F}^{-1} M_m \mathcal{F} = T\) because \(\mathcal{F}^* = \mathcal{F}^{-1}\) and \(m\) is real valued.)

Consequently, for every \(\lambda \in \mathbb{C}\),

\begin{equation*} T - \lambda I = \mathcal{F}^{-1} (M_m - \lambda I) \mathcal{F}, \end{equation*}

so \(T - \lambda I\) is invertible if and only if \(M_m - \lambda I\) is invertible. Hence \(\operatorname{sp}(T) = \operatorname{sp}(M_m)\).

Step 2: \(\operatorname{sp}(M_m) \subseteq [0,1]\).

Suppose \(\lambda \in \mathbb{C} \setminus [0,1]\). Because \([0,1]\) is closed and \(m(t) \in (0,1] \subseteq [0,1]\) for every \(t\), the number

\begin{equation*} d = \inf_{t \in \mathbb{R}} |m(t) - \lambda| \ge \operatorname{dist}(\lambda, [0,1]) \end{equation*}

is positive. Thus \(\frac{1}{m - \lambda}\) is a bounded measurable function on \(\mathbb{R}\) with supremum norm at most \(\frac{1}{d}\), so the multiplication operator \(M_{1/(m - \lambda)}\) is a bounded operator on \(L^2(\mathbb{R})\). Clearly

\begin{equation*} M_{1/(m-\lambda)} (M_m - \lambda I) = (M_m - \lambda I) M_{1/(m - \lambda)} = I , \end{equation*}

since multiplication operators commute and \(\frac{1}{m-\lambda}(m - \lambda) = 1\). Hence \(M_m - \lambda I\) is invertible and \(\lambda \notin \operatorname{sp}(M_m)\).

Step 3: \([0,1] \subseteq \operatorname{sp}(M_m)\).

Let \(\lambda \in [0,1]\). First we produce \(t_1, t_2, \ldots \in \mathbb{R}\) with \(m(t_n) \to \lambda\). If \(\lambda \in (0,1]\), then because \(m\) is continuous with \(m(0) = 1\) and \(m(t) \to 0\) as \(t \to \infty\), the Intermediate Value Theorem provides \(t \in [0, \infty)\) with \(m(t) = \lambda\); take \(t_n = t\) for all \(n\). If \(\lambda = 0\), take \(t_n = n\), so that \(m(t_n) = e^{-2\pi y n} \to 0\).

For each \(n\), the continuity of \(m\) at \(t_n\) gives \(\delta_n > 0\) such that

\begin{equation*} |m(s) - m(t_n)| < \tfrac{1}{n} \quad \text{whenever } |s - t_n| < \delta_n . \end{equation*}

Set \(I_n = (t_n - \delta_n, t_n + \delta_n)\) and

\begin{equation*} h_n = \frac{\chi_{I_n}}{\sqrt{2 \delta_n}} . \end{equation*}

Then \(\|h_n\|_2 = 1\) and

\begin{equation*} \|(M_m - \lambda I) h_n\|_2^2 = \frac{1}{2 \delta_n} \int_{I_n} |m(s) - \lambda|^2\,ds \le \sup_{s \in I_n} |m(s) - \lambda|^2 \le \Bigl( \tfrac{1}{n} + |m(t_n) - \lambda| \Bigr)^2 , \end{equation*}

which tends to \(0\) as \(n \to \infty\).

If \(M_m - \lambda I\) were invertible, with bounded inverse \(S\), then

\begin{equation*} 1 = \|h_n\|_2 = \|S (M_m - \lambda I) h_n \|_2 \le \|S\| \, \|(M_m - \lambda I) h_n\|_2 \to 0 , \end{equation*}

a contradiction. Hence \(M_m - \lambda I\) is not invertible, so \(\lambda \in \operatorname{sp}(M_m)\).

Conclusion. Steps 2 and 3 give \(\operatorname{sp}(M_m) = [0,1]\), and Step 1 gives \(\operatorname{sp}(T) = \operatorname{sp}(M_m)\). Therefore

\begin{equation*} \operatorname{sp}(T) = [0,1] . \end{equation*}

As the book’s remark notes, \([0,1]\) is uncountable, whereas the spectrum of a compact operator is countable by 10.93; hence \(T\) is not compact.

Problem (11C.20)

Prove that if \(f \in L^1(\mathbb{R})\) and \(g \in L^2(\mathbb{R})\), then \(\mathcal{F}(f * g) = \widehat{f} \, \mathcal{F} g\).

Solution

Suppose \(f \in L^1(\mathbb{R})\) and \(g \in L^2(\mathbb{R})\).

Both sides of the asserted equation make sense. On the left, 11.64 (with \(p = 2\)) says that \((f * g)(x)\) is defined for almost every \(x \in \mathbb{R}\) and that

\begin{equation*} \|f * g\|_2 \le \|f\|_1 \|g\|_2 < \infty , \end{equation*}

so \(f * g \in L^2(\mathbb{R})\) and \(\mathcal{F}(f*g)\) is defined by 11.86. On the right, \(\widehat{f} \in L^{\infty}(\mathbb{R})\) with \(\|\widehat{f}\,\|_{\infty} \le \|f\|_1\) by 11.49, and \(\mathcal{F} g \in L^2(\mathbb{R})\), so \(\widehat{f} \, \mathcal{F} g \in L^2(\mathbb{R})\).

The strategy is to verify the equation for \(g\) in the dense subspace \(L^1(\mathbb{R}) \cap L^2(\mathbb{R})\), where it reduces to 11.66, and then pass to the limit.

Step 1: truncation. For \(n \in \mathbb{Z}^+\) define \(g_n = g \chi_{[-n,n]}\). Then \(g_n \in L^2(\mathbb{R})\), and Hölder’s inequality (7.9) gives

\begin{equation*} \|g_n\|_1 = \int_{-n}^{n} |g| \le \sqrt{2n} \, \|g\|_2 < \infty , \end{equation*}

so \(g_n \in L^1(\mathbb{R}) \cap L^2(\mathbb{R})\). Furthermore \(|g - g_n|^2 = |g|^2 \chi_{\mathbb{R} \setminus [-n,n]} \to 0\) pointwise and is dominated by \(|g|^2 \in L^1(\mathbb{R})\), so the Dominated Convergence Theorem (3.31) gives

\begin{equation*} \|g - g_n\|_2 \to 0 \quad \text{as } n \to \infty . \end{equation*}

Step 2: the equation holds for each \(g_n\). Because \(f \in L^1(\mathbb{R})\) and \(g_n \in L^1(\mathbb{R})\), 11.66 gives

\begin{equation*} (f * g_n)\widehat{\phantom{f}} = \widehat{f} \, \widehat{g_n} . \end{equation*}

Moreover \(f * g_n \in L^1(\mathbb{R})\) by 11.64 with \(p = 1\), and \(f * g_n \in L^2(\mathbb{R})\) by 11.64 with \(p = 2\); thus \(f * g_n \in L^1(\mathbb{R}) \cap L^2(\mathbb{R})\), so by the definition of \(\mathcal{F}\) in 11.86 we have \(\mathcal{F}(f * g_n) = (f * g_n)\widehat{\phantom{f}}\). Similarly \(\mathcal{F} g_n = \widehat{g_n}\) because \(g_n \in L^1(\mathbb{R}) \cap L^2(\mathbb{R})\). Hence

\begin{equation*} \mathcal{F}(f * g_n) = \widehat{f} \, \mathcal{F} g_n \qquad \text{for every } n \in \mathbb{Z}^+ . \end{equation*}

Step 3: passing to the limit. For the left side, 11.64 and the linearity of convolution in its second slot give

\begin{equation*} \|f * g - f * g_n\|_2 = \|f * (g - g_n)\|_2 \le \|f\|_1 \|g - g_n\|_2 \to 0 , \end{equation*}

so \(f * g_n \to f * g\) in \(L^2(\mathbb{R})\). Because \(\mathcal{F}\) is a bounded operator on \(L^2(\mathbb{R})\) (indeed an isometry, by 11.82 and 11.86), it follows that

\begin{equation*} \|\mathcal{F}(f * g) - \mathcal{F}(f * g_n)\|_2 = \|\mathcal{F}(f*g - f*g_n)\|_2 = \|f * g - f * g_n\|_2 \to 0 . \end{equation*}

For the right side,

\begin{equation*} \|\widehat{f} \, \mathcal{F} g - \widehat{f} \, \mathcal{F} g_n\|_2 = \|\widehat{f} \, \mathcal{F}(g - g_n)\|_2 \le \|\widehat{f}\,\|_{\infty} \|\mathcal{F}(g - g_n)\|_2 = \|\widehat{f}\,\|_{\infty} \|g - g_n\|_2 \to 0 , \end{equation*}

using again that \(\mathcal{F}\) is an isometry on \(L^2(\mathbb{R})\) and that \(\|\widehat{f}\,\|_{\infty} \le \|f\|_1 < \infty\).

Thus in \(L^2(\mathbb{R})\) we have

\begin{equation*} \mathcal{F}(f * g) = \lim_{n \to \infty} \mathcal{F}(f * g_n) = \lim_{n \to \infty} \widehat{f} \, \mathcal{F} g_n = \widehat{f} \, \mathcal{F} g , \end{equation*}

where the middle equality is Step 2 and the limits are in the norm of \(L^2(\mathbb{R})\) (limits in a normed vector space are unique). Hence \(\mathcal{F}(f * g) = \widehat{f} \, \mathcal{F} g\), meaning that these two functions agree almost everywhere on \(\mathbb{R}\).

Problem (11C.21)

Prove that if \(f, g \in L^2(\mathbb{R})\), then \(\widehat{fg} = (\mathcal{F} f) * (\mathcal{F} g)\).

Solution

Suppose \(f, g \in L^2(\mathbb{R})\).

Both sides make sense pointwise. By Hölder’s inequality (7.9), \(fg \in L^1(\mathbb{R})\) with \(\|fg\|_1 \le \|f\|_2 \|g\|_2\); hence \(\widehat{fg}\) is defined at every \(t \in \mathbb{R}\) by the integral

\begin{equation*} \widehat{fg}(t) = \int_{-\infty}^{\infty} f(x) g(x) e^{-2\pi i t x}\,dx . \end{equation*}

Also \(\mathcal{F} f\) and \(\mathcal{F} g\) are in \(L^2(\mathbb{R})\), so by Hölder’s inequality and the translation invariance of Lebesgue measure the integral defining \(\bigl( (\mathcal{F}f) * (\mathcal{F}g) \bigr)(t)\) converges absolutely for every \(t \in \mathbb{R}\).

We first record two identities for \(\mathcal{F}\) on \(L^2(\mathbb{R})\). Recall that \(L^1(\mathbb{R}) \cap L^2(\mathbb{R})\) is dense in \(L^2(\mathbb{R})\) and that \(\mathcal{F}\) is the unique continuous extension to \(L^2(\mathbb{R})\) of the map \(h \mapsto \widehat{h}\) on \(L^1(\mathbb{R}) \cap L^2(\mathbb{R})\) (see 11.86).

Identity A (conjugation). If \(h \in L^2(\mathbb{R})\), then

\begin{equation*} \mathcal{F}(\overline{h})(s) = \overline{(\mathcal{F}h)(-s)} \quad \text{for almost every } s \in \mathbb{R}. \end{equation*}

Proof: if \(h \in L^1(\mathbb{R}) \cap L^2(\mathbb{R})\), then \(\overline{h} \in L^1(\mathbb{R}) \cap L^2(\mathbb{R})\) and

\begin{equation*} \widehat{\overline{h}}(s) = \int_{-\infty}^{\infty} \overline{h(x)} e^{-2\pi i s x}\,dx = \overline{ \int_{-\infty}^{\infty} h(x) e^{2 \pi i s x}\,dx } = \overline{\widehat{h}(-s)} . \end{equation*}

Thus the two maps \(h \mapsto \mathcal{F}(\overline{h})\) and \(h \mapsto \overline{(\mathcal{F}h)(-\,\cdot\,)}\) from \(L^2(\mathbb{R})\) to \(L^2(\mathbb{R})\) agree on \(L^1(\mathbb{R}) \cap L^2(\mathbb{R})\). Each of these maps is conjugate linear and norm preserving (because \(\mathcal{F}\) is an isometry by 11.82, and because conjugation and reflection preserve the \(L^2\) norm), hence each is continuous. Two continuous maps that agree on a dense subset agree everywhere, proving Identity A.

Identity B (modulation). If \(h \in L^2(\mathbb{R})\) and \(b \in \mathbb{R}\) and \(w(x) = e^{2 \pi i b x} h(x)\), then

\begin{equation*} (\mathcal{F} w)(s) = (\mathcal{F} h)(s - b) \quad \text{for almost every } s \in \mathbb{R}. \end{equation*}

Proof: for \(h \in L^1(\mathbb{R}) \cap L^2(\mathbb{R})\) this is exactly 11.55(b) (note \(|w| = |h|\), so \(w\) is also in \(L^1(\mathbb{R}) \cap L^2(\mathbb{R})\)). Both \(h \mapsto \mathcal{F}(e^{2\pi i b \,\cdot\,} h)\) and \(h \mapsto (\mathcal{F}h)(\,\cdot\, - b)\) are linear norm-preserving maps of \(L^2(\mathbb{R})\) into itself, hence continuous, and they agree on the dense subspace \(L^1(\mathbb{R}) \cap L^2(\mathbb{R})\); therefore they agree on \(L^2(\mathbb{R})\).

Now fix \(t \in \mathbb{R}\) and define \(h_t : \mathbb{R} \to \mathbb{C}\) by

\begin{equation*} h_t(x) = \overline{g(x)}\, e^{2 \pi i t x} . \end{equation*}

Then \(|h_t| = |g|\), so \(h_t \in L^2(\mathbb{R})\), and

\begin{equation*} \langle f, h_t \rangle = \int_{-\infty}^{\infty} f(x) \overline{h_t(x)}\,dx = \int_{-\infty}^{\infty} f(x) g(x) e^{-2 \pi i t x}\,dx = \widehat{fg}(t) . \end{equation*}

By 11.87(a) the operator \(\mathcal{F}\) is unitary on \(L^2(\mathbb{R})\), so it preserves inner products:

\begin{equation*} \widehat{fg}(t) = \langle f, h_t \rangle = \langle \mathcal{F}f, \mathcal{F}h_t \rangle . \end{equation*}

Next we identify \(\mathcal{F} h_t\). Applying Identity B with \(b = t\) and \(h = \overline{g}\), and then Identity A with \(h = g\), we get for almost every \(s \in \mathbb{R}\)

\begin{equation*} (\mathcal{F} h_t)(s) = \bigl( \mathcal{F} \overline{g} \bigr)(s - t) = \overline{(\mathcal{F}g)\bigl(-(s-t)\bigr)} = \overline{(\mathcal{F}g)(t - s)} . \end{equation*}

Therefore

\begin{equation*} \widehat{fg}(t) = \langle \mathcal{F}f, \mathcal{F}h_t \rangle = \int_{-\infty}^{\infty} (\mathcal{F}f)(s) \, \overline{(\mathcal{F}h_t)(s)}\,ds = \int_{-\infty}^{\infty} (\mathcal{F}f)(s) \, (\mathcal{F}g)(t - s)\,ds = \bigl( (\mathcal{F}f) * (\mathcal{F}g) \bigr)(t) , \end{equation*}

where changing \(\mathcal{F}h_t\) on a set of measure zero does not affect the value of the integral.

Because \(t \in \mathbb{R}\) was arbitrary, we conclude that

\begin{equation*} \widehat{fg} = (\mathcal{F}f) * (\mathcal{F}g) , \end{equation*}

with equality at every point of \(\mathbb{R}\) (both sides are in fact uniformly continuous functions on \(\mathbb{R}\): the left side by 11.49, the right side by Exercise 10 in this section applied with \(p = 2\)).

Probability Measures

Exercises

Problem (12.1)

Suppose \((\Omega, \mathcal{F}, P)\) is a probability space and \(A \in \mathcal{F}\). Prove that \(A\) and \(\Omega \setminus A\) are independent if and only if \(P(A) = 0\) or \(P(A) = 1\).

Solution

Because \(A \cap (\Omega \setminus A) = \emptyset\), we have \(P\bigl(A \cap (\Omega \setminus A)\bigr) = P(\emptyset) = 0\). Also, \(P(\Omega \setminus A) = 1 - P(A)\), because \(P(\Omega) = 1\) and \(P\) is finitely additive.

Thus, by the definition of independent events (12.7), the events \(A\) and \(\Omega \setminus A\) are independent if and only if

\begin{equation*} 0 = P(A) \cdot \bigl(1 - P(A)\bigr). \end{equation*}

Since \(P(A)\) is a real number, the product \(P(A)\bigl(1 - P(A)\bigr)\) equals \(0\) if and only if \(P(A) = 0\) or \(1 - P(A) = 0\), in other words if and only if \(P(A) = 0\) or \(P(A) = 1\).

Hence \(A\) and \(\Omega \setminus A\) are independent if and only if \(P(A) = 0\) or \(P(A) = 1\), as desired.

Problem (12.2)

Suppose \(P\) is Lebesgue measure on \([0, 1]\). Give an example of two disjoint Borel subsets \(A\) and \(B\) of \([0, 1]\) such that \(P(A) = P(B) = \frac{1}{2}\), \([0, \frac{1}{2}]\) and \(A\) are independent, and \([0, \frac{1}{2}]\) and \(B\) are independent.

Solution

Take

\begin{equation*} A = \left[0, \tfrac{1}{4}\right] \cup \left(\tfrac{1}{2}, \tfrac{3}{4}\right], \qquad B = \left(\tfrac{1}{4}, \tfrac{1}{2}\right] \cup \left(\tfrac{3}{4}, 1\right]. \end{equation*}

Each of \(A\) and \(B\) is a union of two intervals and hence is a Borel subset of \([0,1]\). The four intervals

\begin{equation*} \left[0, \tfrac14\right], \quad \left(\tfrac14, \tfrac12\right], \quad \left(\tfrac12, \tfrac34\right], \quad \left(\tfrac34, 1\right] \end{equation*}

are pairwise disjoint and their union is \([0,1]\); since \(A\) is the union of the first and third while \(B\) is the union of the second and fourth, the sets \(A\) and \(B\) are disjoint.

Lebesgue measure of a closed interval equals its length (2.14, together with 2.68, which says that outer measure is a measure on the Borel sets), and a one-point set has measure \(0\); hence any interval with endpoints \(a < b\) has measure \(b - a\), regardless of whether the endpoints are included. So each of the four intervals above has measure \(\frac14\). By additivity,

\begin{equation*} P(A) = \tfrac14 + \tfrac14 = \tfrac12, \qquad P(B) = \tfrac14 + \tfrac14 = \tfrac12 . \end{equation*}

Also \(P\bigl(\left[0, \tfrac12\right]\bigr) = \tfrac12\).

Now

\begin{equation*} \left[0, \tfrac12\right] \cap A = \left[0, \tfrac14\right], \end{equation*}

because the part \(\left(\tfrac12, \tfrac34\right]\) of \(A\) is disjoint from \(\left[0,\tfrac12\right]\). Hence

\begin{equation*} P\left(\left[0, \tfrac12\right] \cap A\right) = \tfrac14 = \tfrac12 \cdot \tfrac12 = P\left(\left[0, \tfrac12\right]\right) \cdot P(A), \end{equation*}

so \(\left[0, \tfrac12\right]\) and \(A\) are independent events.

Similarly,

\begin{equation*} \left[0, \tfrac12\right] \cap B = \left(\tfrac14, \tfrac12\right], \end{equation*}

and therefore

\begin{equation*} P\left(\left[0, \tfrac12\right] \cap B\right) = \tfrac14 = \tfrac12 \cdot \tfrac12 = P\left(\left[0, \tfrac12\right]\right) \cdot P(B), \end{equation*}

so \(\left[0, \tfrac12\right]\) and \(B\) are independent events.

Thus \(A\) and \(B\) have all the required properties.

Problem (12.3)

Suppose \((\Omega, \mathcal{F}, P)\) is a probability space and \(A, B \in \mathcal{F}\). Prove that the following are equivalent.

  • \(A\) and \(B\) are independent events.
  • \(A\) and \(\Omega \setminus B\) are independent events.
  • \(\Omega \setminus A\) and \(B\) are independent events.
  • \(\Omega \setminus A\) and \(\Omega \setminus B\) are independent events.
Solution

We first prove a lemma from which everything follows.

Lemma: if \(C, D \in \mathcal{F}\), then \(C\) and \(D\) are independent if and only if \(\Omega \setminus C\) and \(D\) are independent.

To prove the lemma, note that \(D\) is the disjoint union of \(C \cap D\) and \((\Omega \setminus C) \cap D\). Hence, by additivity of \(P\),

\begin{equation*} P\bigl((\Omega \setminus C) \cap D\bigr) = P(D) - P(C \cap D). \end{equation*}

Also, since \(P(\Omega \setminus C) = 1 - P( C)\),

\begin{equation*} P(\Omega \setminus C) \cdot P(D) = P(D) - P( C) \cdot P(D). \end{equation*}

Subtracting the second displayed equation from the first, we obtain

\begin{equation*} P\bigl((\Omega \setminus C) \cap D\bigr) - P(\Omega \setminus C) \cdot P(D) = P( C) \cdot P(D) - P(C \cap D). \end{equation*}

Thus the left side equals \(0\) if and only if the right side equals \(0\); in other words, \(\Omega \setminus C\) and \(D\) are independent if and only if \(C\) and \(D\) are independent. This proves the lemma.

Note also that independence is a symmetric relation: \(P(C \cap D) = P( C) \cdot P(D)\) if and only if \(P(D \cap C) = P(D) \cdot P( C)\), because \(C \cap D = D \cap C\) and multiplication of real numbers is commutative. Hence the lemma may equally be applied to complement the second event.

Now we prove the equivalences.

Applying the lemma with \(C = B\) and \(D = A\) (and using symmetry of independence) shows that the first bullet point is equivalent to the second bullet point: indeed \(A\) and \(B\) are independent if and only if \(B\) and \(A\) are independent, which by the lemma holds if and only if \(\Omega \setminus B\) and \(A\) are independent, which by symmetry holds if and only if \(A\) and \(\Omega \setminus B\) are independent.

Applying the lemma with \(C = A\) and \(D = B\) shows that the first bullet point (\(A\) and \(B\) independent) is equivalent to the third bullet point (\(\Omega \setminus A\) and \(B\) independent).

Finally, applying the lemma with \(C = A\) and \(D = \Omega \setminus B\) shows that the second bullet point (\(A\) and \(\Omega \setminus B\) independent) is equivalent to the fourth bullet point (\(\Omega \setminus A\) and \(\Omega \setminus B\) independent).

Thus all four statements are equivalent to the first one, and hence they are all equivalent to each other.

Problem (12.4)

Suppose \((\Omega, \mathcal{F}, P)\) is a probability space and \(\{A_k\}_{k \in \Gamma}\) is a family of events. Prove the family \(\{A_k\}_{k \in \Gamma}\) is independent if and only if the family \(\{\Omega \setminus A_k\}_{k \in \Gamma}\) is independent.

Solution

Suppose first that \(\{A_k\}_{k \in \Gamma}\) is independent. We prove the following stronger statement, which allows an arbitrary mixture of the events and their complements.

Claim: for every integer \(n \ge 0\), every \(m \in \{0, 1, \ldots, n\}\), and all distinct \(k_1, \ldots, k_n \in \Gamma\),

\begin{equation*} P\Bigl(\bigcap_{j=1}^{m} (\Omega \setminus A_{k_j}) \;\cap \bigcap_{j=m+1}^{n} A_{k_j}\Bigr) = \prod_{j=1}^{m}\bigl(1 - P(A_{k_j})\bigr) \cdot \prod_{j=m+1}^{n} P(A_{k_j}), \end{equation*}

where an empty intersection is interpreted as \(\Omega\) and an empty product is interpreted as \(1\).

We prove the claim by induction on \(m\) (the statement for each fixed \(m\) being asserted for all \(n \ge m\) and all choices of distinct indices).

For the base case \(m = 0\) the claim reads

\begin{equation*} P(A_{k_1} \cap \cdots \cap A_{k_n}) = P(A_{k_1}) \cdots P(A_{k_n}), \end{equation*}

which for \(n \ge 1\) is exactly the hypothesis that \(\{A_k\}_{k \in \Gamma}\) is independent (12.7), and which for \(n = 0\) reads \(P(\Omega) = 1\), true because \(P\) is a probability measure.

Now suppose \(m \ge 1\) and the claim holds with \(m\) replaced by \(m - 1\). Let \(n \ge m\) and let \(k_1, \ldots, k_n \in \Gamma\) be distinct. Set

\begin{equation*} C = \bigcap_{j=1}^{m-1} (\Omega \setminus A_{k_j}) \;\cap \bigcap_{j=m+1}^{n} A_{k_j}, \end{equation*}

which is the intersection corresponding to the \(n - 1\) distinct indices \(k_1, \ldots, k_{m-1}, k_{m+1}, \ldots, k_n\), with the first \(m - 1\) of them complemented. Hence the induction hypothesis applies and gives

\begin{equation*} P( C) = \prod_{j=1}^{m-1}\bigl(1 - P(A_{k_j})\bigr) \cdot \prod_{j=m+1}^{n} P(A_{k_j}). \end{equation*}

Also, \(C \cap A_{k_m}\) is the intersection corresponding to the \(n\) distinct indices \(k_1, \ldots, k_{m-1}, k_m, k_{m+1}, \ldots, k_n\) with exactly the first \(m-1\) complemented (reorder the list so that the complemented indices come first, which does not change the intersection or the product); thus the induction hypothesis also gives

\begin{equation*} P(C \cap A_{k_m}) = \prod_{j=1}^{m-1}\bigl(1 - P(A_{k_j})\bigr) \cdot P(A_{k_m}) \cdot \prod_{j=m+1}^{n} P(A_{k_j}). \end{equation*}

Because \(C\) is the disjoint union of \(C \cap A_{k_m}\) and \(C \cap (\Omega \setminus A_{k_m})\), additivity of \(P\) gives

\begin{equation*} \begin{aligned} P\bigl(C \cap (\Omega \setminus A_{k_m})\bigr) &= P( C) - P(C \cap A_{k_m}) \\ &= \prod_{j=1}^{m-1}\bigl(1 - P(A_{k_j})\bigr) \cdot \bigl(1 - P(A_{k_m})\bigr) \cdot \prod_{j=m+1}^{n} P(A_{k_j}) \\ &= \prod_{j=1}^{m}\bigl(1 - P(A_{k_j})\bigr) \cdot \prod_{j=m+1}^{n} P(A_{k_j}). \end{aligned} \end{equation*}

Since \(C \cap (\Omega \setminus A_{k_m}) = \bigcap_{j=1}^{m} (\Omega \setminus A_{k_j}) \cap \bigcap_{j=m+1}^{n} A_{k_j}\), this completes the induction and proves the claim.

Taking \(m = n\) in the claim, we obtain

\begin{equation*} P\bigl((\Omega \setminus A_{k_1}) \cap \cdots \cap (\Omega \setminus A_{k_n})\bigr) = \bigl(1 - P(A_{k_1})\bigr) \cdots \bigl(1 - P(A_{k_n})\bigr) = P(\Omega \setminus A_{k_1}) \cdots P(\Omega \setminus A_{k_n}) \end{equation*}

for all distinct \(k_1, \ldots, k_n \in \Gamma\), where the last equality holds because \(P(\Omega \setminus A_k) = 1 - P(A_k)\). Thus \(\{\Omega \setminus A_k\}_{k \in \Gamma}\) is an independent family of events.

For the other direction, suppose \(\{\Omega \setminus A_k\}_{k \in \Gamma}\) is independent. Applying the implication just proved to the family \(\{\Omega \setminus A_k\}_{k \in \Gamma}\) shows that \(\{\Omega \setminus (\Omega \setminus A_k)\}_{k \in \Gamma}\) is independent. Because \(\Omega \setminus (\Omega \setminus A_k) = A_k\) for each \(k \in \Gamma\), this says that \(\{A_k\}_{k \in \Gamma}\) is independent, completing the proof.

Problem (12.5)

Give an example of a probability space \((\Omega, \mathcal{F}, P)\) and events \(A, B_1, B_2\) such that \(A\) and \(B_1\) are independent, \(A\) and \(B_2\) are independent, but \(A\) and \(B_1 \cup B_2\) are not independent.

Solution

Let \(\Omega = \{1, 2, 3, 4\}\), let \(\mathcal{F}\) be the collection of all subsets of \(\Omega\), and let \(P\) be counting measure on \(\Omega\) divided by \(4\); thus \(P(E) = \frac{|E|}{4}\) for each \(E \subseteq \Omega\), where \(|E|\) denotes the number of elements of \(E\). Then \(P(\Omega) = 1\), so \((\Omega, \mathcal{F}, P)\) is a probability space.

Define

\begin{equation*} A = \{1, 2\}, \qquad B_1 = \{1, 3\}, \qquad B_2 = \{1, 4\}. \end{equation*}

Then \(P(A) = P(B_1) = P(B_2) = \frac{1}{2}\).

Because \(A \cap B_1 = \{1\}\), we have

\begin{equation*} P(A \cap B_1) = \tfrac14 = \tfrac12 \cdot \tfrac12 = P(A) \cdot P(B_1), \end{equation*}

so \(A\) and \(B_1\) are independent events.

Because \(A \cap B_2 = \{1\}\), we similarly have

\begin{equation*} P(A \cap B_2) = \tfrac14 = \tfrac12 \cdot \tfrac12 = P(A) \cdot P(B_2), \end{equation*}

so \(A\) and \(B_2\) are independent events.

However, \(B_1 \cup B_2 = \{1, 3, 4\}\), so \(P(B_1 \cup B_2) = \frac{3}{4}\), while \(A \cap (B_1 \cup B_2) = \{1\}\), so \(P\bigl(A \cap (B_1 \cup B_2)\bigr) = \frac{1}{4}\). Since

\begin{equation*} \tfrac14 \ne \tfrac12 \cdot \tfrac34 = \tfrac38, \end{equation*}

the events \(A\) and \(B_1 \cup B_2\) are not independent.

The reason the example works is visible from inclusion-exclusion: independence of \(A\) with \(B_1\) and with \(B_2\) gives

\begin{equation*} P\bigl(A \cap (B_1 \cup B_2)\bigr) - P(A)P(B_1 \cup B_2) = P(A)P(B_1 \cap B_2) - P(A \cap B_1 \cap B_2), \end{equation*}

so \(A\) and \(B_1 \cup B_2\) are independent precisely when \(A\) and \(B_1 \cap B_2\) are independent. In the example above \(B_1 \cap B_2 = \{1\}\) has probability \(\frac14\) while \(A \cap B_1 \cap B_2 = \{1\}\) has probability \(\frac14 \ne \frac12 \cdot \frac14\).

Problem (12.6)

Give an example of a probability space \((\Omega, \mathcal{F}, P)\) and events \(A_1, A_2, A_3\) such that \(A_1\) and \(A_2\) are independent, \(A_1\) and \(A_3\) are independent, and \(A_2\) and \(A_3\) are independent, but the family \(A_1, A_2, A_3\) is not independent.

Solution

Let \(\Omega = \{H, T\}^2\) be the sample space of two tosses of a fair coin, so \(\Omega\) has the four elements \((H,H)\), \((H,T)\), \((T,H)\), \((T,T)\). Let \(\mathcal{F}\) be the collection of all subsets of \(\Omega\) and let \(P\) be counting measure on \(\Omega\) divided by \(4\). Then \(P(\Omega) = 1\), so \((\Omega, \mathcal{F}, P)\) is a probability space.

Define

\begin{equation*} \begin{aligned} A_1 &= \{\omega \in \Omega : \omega_1 = H\} = \{(H,H), (H,T)\}, \\ A_2 &= \{\omega \in \Omega : \omega_2 = H\} = \{(H,H), (T,H)\}, \\ A_3 &= \{\omega \in \Omega : \omega_1 = \omega_2\} = \{(H,H), (T,T)\}. \end{aligned} \end{equation*}

Thus \(A_1\) is the event that the first toss is heads, \(A_2\) is the event that the second toss is heads, and \(A_3\) is the event that the two tosses agree. Each of these three sets has two elements, so

\begin{equation*} P(A_1) = P(A_2) = P(A_3) = \tfrac12 . \end{equation*}

Each of the three pairwise intersections equals \(\{(H,H)\}\):

\begin{equation*} A_1 \cap A_2 = A_1 \cap A_3 = A_2 \cap A_3 = \{(H,H)\}. \end{equation*}

Indeed, \(A_1 \cap A_2\) consists of the outcomes with both tosses heads; \(A_1 \cap A_3\) consists of the outcomes whose first toss is heads and whose tosses agree, hence both heads; and \(A_2 \cap A_3\) consists of the outcomes whose second toss is heads and whose tosses agree, hence both heads.

Therefore each pairwise intersection has probability \(\frac14\), and

\begin{equation*} P(A_j \cap A_k) = \tfrac14 = \tfrac12 \cdot \tfrac12 = P(A_j) \cdot P(A_k) \end{equation*}

whenever \(j, k \in \{1, 2, 3\}\) with \(j \ne k\). Hence \(A_1\) and \(A_2\) are independent, \(A_1\) and \(A_3\) are independent, and \(A_2\) and \(A_3\) are independent.

However,

\begin{equation*} A_1 \cap A_2 \cap A_3 = \{(H,H)\}, \end{equation*}

so

\begin{equation*} P(A_1 \cap A_2 \cap A_3) = \tfrac14 \ne \tfrac18 = P(A_1) \cdot P(A_2) \cdot P(A_3). \end{equation*}

Thus the family \(A_1, A_2, A_3\) fails the condition in the definition of an independent family of events (12.7) for the index set \(\{1, 2, 3\}\), and hence \(A_1, A_2, A_3\) is not an independent family.

Problem (12.7)

Suppose \((\Omega, \mathcal{F}, P)\) is a probability space, \(A \in \mathcal{F}\), and \(B_1 \subseteq B_2 \subseteq \cdots\) is an increasing sequence of events such that \(A\) and \(B_n\) are independent events for each \(n \in \mathbf{Z}^+\). Show that \(A\) and \(\bigcup_{n=1}^{\infty} B_n\) are independent.

Solution

Let \(B = \bigcup_{n=1}^{\infty} B_n\). Because \(\mathcal{F}\) is a \(\sigma\)-algebra and each \(B_n \in \mathcal{F}\), we have \(B \in \mathcal{F}\); also \(A \cap B_n \in \mathcal{F}\) for each \(n \in \mathbf{Z}^+\).

The sequence \(B_1 \subseteq B_2 \subseteq \cdots\) is increasing, so the sequence

\begin{equation*} A \cap B_1 \subseteq A \cap B_2 \subseteq \cdots \end{equation*}

is also increasing. Furthermore, distributing the intersection over the union gives

\begin{equation*} \bigcup_{n=1}^{\infty} (A \cap B_n) = A \cap \bigcup_{n=1}^{\infty} B_n = A \cap B . \end{equation*}

Now apply the formula for the measure of an increasing union (2.59) twice, once to the sequence \(\{A \cap B_n\}_{n \in \mathbf{Z}^+}\) and once to the sequence \(\{B_n\}_{n \in \mathbf{Z}^+}\):

\begin{equation*} P(A \cap B) = P\Bigl(\bigcup_{n=1}^{\infty} (A \cap B_n)\Bigr) = \lim_{n \to \infty} P(A \cap B_n) \end{equation*}

and

\begin{equation*} P(B) = \lim_{n \to \infty} P(B_n). \end{equation*}

By hypothesis, \(A\) and \(B_n\) are independent events for each \(n \in \mathbf{Z}^+\), which means \(P(A \cap B_n) = P(A) \cdot P(B_n)\). Hence

\begin{equation*} P(A \cap B) = \lim_{n \to \infty} P(A \cap B_n) = \lim_{n \to \infty} P(A) \cdot P(B_n) = P(A) \cdot \lim_{n \to \infty} P(B_n) = P(A) \cdot P(B), \end{equation*}

where the third equality holds because \(P(A)\) is a constant and the limit \(\lim_{n \to \infty} P(B_n)\) exists (it equals \(P(B)\), a real number, by the previous paragraph).

Thus \(P(A \cap B) = P(A) \cdot P(B)\), which by 12.7 means that \(A\) and \(B = \bigcup_{n=1}^{\infty} B_n\) are independent events, as desired.

Problem (12.8)

Suppose \((\Omega, \mathcal{F}, P)\) is a probability space and \(\{A_t\}_{t \in \mathbb{R}}\) is an independent family of events such that \(P(A_t) < 1\) for each \(t \in \mathbb{R}\). Prove that there exists a sequence \(t_1, t_2, \ldots\) in \(\mathbb{R}\) such that

\begin{equation*} P\Bigl( \bigcap_{n=1}^{\infty} A_{t_n} \Bigr) = 0 . \end{equation*}

Solution

For each \(t \in \mathbb{R}\) we have \(P(\Omega \setminus A_t) = 1 - P(A_t) > 0\), because \(P(\Omega) = 1\) and \(P\) is finitely additive. For \(m \in \mathbb{Z}^+\) let

\begin{equation*} S_m = \Bigl\{ t \in \mathbb{R} : P(\Omega \setminus A_t) > \tfrac{1}{m} \Bigr\} . \end{equation*}

Since \(P(\Omega \setminus A_t) > 0\) for every \(t\), every \(t \in \mathbb{R}\) lies in \(S_m\) for all sufficiently large \(m\); hence

\begin{equation*} \mathbb{R} = \bigcup_{m=1}^{\infty} S_m . \end{equation*}

If every \(S_m\) were finite, then \(\mathbb{R}\) would be a countable union of finite sets and therefore countable, which is false. Thus there exists \(m \in \mathbb{Z}^+\) such that \(S_m\) is infinite. Fix such an \(m\) and choose distinct elements

\begin{equation*} t_1, t_2, t_3, \ldots \in S_m . \end{equation*}

Fix \(N \in \mathbb{Z}^+\). Because \(t_1, \ldots, t_N\) are distinct elements of the index set \(\mathbb{R}\) and the family \(\{A_t\}_{t \in \mathbb{R}}\) is independent, the definition of independence (12.7) gives

\begin{equation*} P\Bigl( \bigcap_{n=1}^{N} A_{t_n} \Bigr) = \prod_{n=1}^{N} P(A_{t_n}) . \end{equation*}

For each \(n\) we have \(t_n \in S_m\), so

\begin{equation*} P(A_{t_n}) = 1 - P(\Omega \setminus A_{t_n}) < 1 - \tfrac{1}{m} . \end{equation*}

Therefore

\begin{equation*} P\Bigl( \bigcap_{n=1}^{N} A_{t_n} \Bigr) < \Bigl( 1 - \tfrac{1}{m} \Bigr)^{N} . \end{equation*}

Now \(\bigcap_{n=1}^{\infty} A_{t_n} \subseteq \bigcap_{n=1}^{N} A_{t_n}\), so monotonicity of the measure \(P\) (2.57) gives

\begin{equation*} 0 \le P\Bigl( \bigcap_{n=1}^{\infty} A_{t_n} \Bigr) \le \Bigl( 1 - \tfrac{1}{m} \Bigr)^{N} \end{equation*}

for every \(N \in \mathbb{Z}^+\). Since \(0 \le 1 - \frac{1}{m} < 1\), the right side has limit \(0\) as \(N \to \infty\). Hence

\begin{equation*} P\Bigl( \bigcap_{n=1}^{\infty} A_{t_n} \Bigr) = 0, \end{equation*}

as desired.

Note where the hypotheses were used: the uncountability of the index set \(\mathbb{R}\) is what forced some \(S_m\) to be infinite, which in turn produced a uniform bound \(P(A_{t_n}) < 1 - \frac{1}{m}\) along the chosen sequence. Without such a uniform bound the conclusion can fail: if \(P(A_n) = 1 - 2^{-n}\) for an independent sequence, then \(P(\bigcap_{n=1}^{\infty} A_n) = \prod_{n=1}^{\infty}(1 - 2^{-n}) > 0\).

Problem (12.9)

Suppose \((\Omega, \mathcal{F}, P)\) is a probability space and \(B_1, \ldots, B_n \in \mathcal{F}\) are such that \(P(B_1 \cap \cdots \cap B_n) > 0\). Prove that

\begin{equation*} P(A \cap B_1 \cap \cdots \cap B_n) = P(B_1) \cdot P_{B_1}(B_2) \cdots P_{B_1 \cap \cdots \cap B_{n-1}}(B_n) \cdot P_{B_1 \cap \cdots \cap B_n}(A) \end{equation*}

for every event \(A \in \mathcal{F}\).

Solution

For \(k \in \{1, \ldots, n\}\) set

\begin{equation*} C_k = B_1 \cap \cdots \cap B_k . \end{equation*}

Then \(C_1 \supseteq C_2 \supseteq \cdots \supseteq C_n\), and \(C_n = B_1 \cap \cdots \cap B_n\) has positive probability by hypothesis. Because \(P\) preserves order (2.57(a)),

\begin{equation*} P(C_k) \ge P(C_n) > 0 \end{equation*}

for each \(k \in \{1, \ldots, n\}\). Hence each of the conditional probability measures \(P_{C_1}, \ldots, P_{C_n}\) is defined (12.23); this is exactly what is needed for the right side of the displayed equation to make sense.

Now fix \(A \in \mathcal{F}\). By the definition of conditional probability (12.23), for each \(k \in \{2, \ldots, n\}\) we have

\begin{equation*} P_{C_{k-1}}(B_k) = \frac{P(B_k \cap C_{k-1})}{P(C_{k-1})} = \frac{P(C_k)}{P(C_{k-1})}, \end{equation*}

where the second equality holds because \(B_k \cap C_{k-1} = B_k \cap B_1 \cap \cdots \cap B_{k-1} = C_k\). Again by 12.23,

\begin{equation*} P_{C_n}(A) = \frac{P(A \cap C_n)}{P(C_n)} . \end{equation*}

All the denominators above are nonzero, so we may multiply these \(n\) equations together, along with the trivial identity \(P(B_1) = P(C_1)\):

\begin{equation*} \begin{aligned} P(B_1) \cdot P_{C_1}(B_2) \cdots P_{C_{n-1}}(B_n) \cdot P_{C_n}(A) &= P(C_1) \cdot \frac{P(C_2)}{P(C_1)} \cdot \frac{P(C_3)}{P(C_2)} \cdots \frac{P(C_n)}{P(C_{n-1})} \cdot \frac{P(A \cap C_n)}{P(C_n)} \\ &= P(A \cap C_n), \end{aligned} \end{equation*}

because the product telescopes: every factor \(P(C_k)\) for \(k \in \{1, \ldots, n\}\) occurs once in a numerator and once in a denominator. Since \(C_n = B_1 \cap \cdots \cap B_n\), the right side equals \(P(A \cap B_1 \cap \cdots \cap B_n)\), which completes the proof.

In the degenerate case \(n = 1\) the displayed formula reads \(P(A \cap B_1) = P(B_1) \cdot P_{B_1}(A)\), which is just the definition 12.23 rearranged, consistent with the above.

Problem (12.10)

Suppose \((\Omega, \mathcal{F}, P)\) is a probability space and \(A \in \mathcal{F}\) is an event such that \(0 < P(A) < 1\). Prove that

\begin{equation*} P(B) = P_A(B) \cdot P(A) + P_{\Omega \setminus A}(B) \cdot P(\Omega \setminus A) \end{equation*}

for every event \(B \in \mathcal{F}\).

Solution

First note that both conditional probability measures appearing in the statement are defined. We have \(P(A) > 0\) by hypothesis, and because \(P(\Omega) = 1\) and \(A\) and \(\Omega \setminus A\) are disjoint events whose union is \(\Omega\), additivity of \(P\) gives

\begin{equation*} P(\Omega \setminus A) = 1 - P(A) > 0, \end{equation*}

the last inequality holding because \(P(A) < 1\). Thus \(P_A\) and \(P_{\Omega \setminus A}\) both make sense (12.23).

Now suppose \(B \in \mathcal{F}\). By the definition of conditional probability (12.23),

\begin{equation*} P_A(B) \cdot P(A) = \frac{P(B \cap A)}{P(A)} \cdot P(A) = P(B \cap A) \end{equation*}

and

\begin{equation*} P_{\Omega \setminus A}(B) \cdot P(\Omega \setminus A) = \frac{P\bigl( B \cap (\Omega \setminus A) \bigr)}{P(\Omega \setminus A)} \cdot P(\Omega \setminus A) = P\bigl( B \cap (\Omega \setminus A) \bigr) . \end{equation*}

The sets \(B \cap A\) and \(B \cap (\Omega \setminus A)\) are disjoint elements of \(\mathcal{F}\) whose union is \(B\), because

\begin{equation*} (B \cap A) \cup \bigl( B \cap (\Omega \setminus A) \bigr) = B \cap \bigl( A \cup (\Omega \setminus A) \bigr) = B \cap \Omega = B . \end{equation*}

Hence countable additivity of the measure \(P\) (applied to the sequence \(B \cap A\), \(B \cap (\Omega \setminus A)\), \(\emptyset\), \(\emptyset\), \(\ldots\)) gives

\begin{equation*} P(B) = P(B \cap A) + P\bigl( B \cap (\Omega \setminus A) \bigr) . \end{equation*}

Combining the three displayed equations above yields

\begin{equation*} P(B) = P_A(B) \cdot P(A) + P_{\Omega \setminus A}(B) \cdot P(\Omega \setminus A), \end{equation*}

as desired. (This is the case \(n = 2\) of the identity 12.26 that appears inside the proof of the second version of Bayes’ Theorem, with \(A_1 = A\) and \(A_2 = \Omega \setminus A\).)

Problem (12.11)

Give an example of a probability space \((\Omega, \mathcal{F}, P)\) and \(X, Y \in \mathcal{L}^2(P)\) such that \(\sigma^2(X + Y) = \sigma^2(X) + \sigma^2(Y)\) but \(X\) and \(Y\) are not independent random variables.

Solution

Let \(\Omega = \{-1, 0, 1\}\), let \(\mathcal{F}\) be the collection of all subsets of \(\Omega\), and let

\begin{equation*} P = \frac{\text{counting measure on } \Omega}{3}, \end{equation*}

which is a probability measure on \((\Omega, \mathcal{F})\) by the first bullet point of 12.2. Define random variables \(X\) and \(Y\) on \(\Omega\) by

\begin{equation*} X(\omega) = \omega \quad \text{and} \quad Y(\omega) = \omega^2 . \end{equation*}

Every function on the finite set \(\Omega\) is \(\mathcal{F}\)-measurable and bounded, so \(X, Y \in \mathcal{L}^2(P)\).

First we compute the relevant expectations. Since \(P(\{\omega\}) = \frac{1}{3}\) for each \(\omega \in \Omega\), the integral over \(\Omega\) of a function \(f\) is \(\frac{1}{3}\bigl( f(-1) + f(0) + f(1) \bigr)\). Hence

\begin{equation*} EX = \tfrac{1}{3}(-1 + 0 + 1) = 0, \qquad E(X^2) = \tfrac{1}{3}(1 + 0 + 1) = \tfrac{2}{3}, \end{equation*}

\begin{equation*} EY = \tfrac{1}{3}(1 + 0 + 1) = \tfrac{2}{3}, \qquad E(Y^2) = \tfrac{1}{3}(1 + 0 + 1) = \tfrac{2}{3}, \end{equation*}

where \(E(Y^2)\) uses \(\omega^4 = \omega^2\) for \(\omega \in \{-1,0,1\}\). Thus the variance formula 12.20 gives

\begin{equation*} \sigma^2(X) = E(X^2) - (EX)^2 = \tfrac{2}{3} - 0 = \tfrac{2}{3}, \end{equation*}

\begin{equation*} \sigma^2(Y) = E(Y^2) - (EY)^2 = \tfrac{2}{3} - \tfrac{4}{9} = \tfrac{2}{9} . \end{equation*}

Hence

\begin{equation*} \sigma^2(X) + \sigma^2(Y) = \tfrac{2}{3} + \tfrac{2}{9} = \tfrac{8}{9} . \end{equation*}

Now consider \(X + Y\). We have \((X+Y)(-1) = -1 + 1 = 0\), \((X+Y)(0) = 0\), and \((X+Y)(1) = 1 + 1 = 2\). Therefore

\begin{equation*} E(X + Y) = \tfrac{1}{3}(0 + 0 + 2) = \tfrac{2}{3}, \qquad E\bigl( (X+Y)^2 \bigr) = \tfrac{1}{3}(0 + 0 + 4) = \tfrac{4}{3}, \end{equation*}

and so by 12.20 again

\begin{equation*} \sigma^2(X + Y) = \tfrac{4}{3} - \tfrac{4}{9} = \tfrac{8}{9} = \sigma^2(X) + \sigma^2(Y) . \end{equation*}

Finally, \(X\) and \(Y\) are not independent random variables. Take the Borel sets \(U = V = \{0\}\). Then

\begin{equation*} \{X \in U\} = \{\omega \in \Omega : \omega = 0\} = \{0\}, \qquad \{Y \in V\} = \{\omega \in \Omega : \omega^2 = 0\} = \{0\}, \end{equation*}

so

\begin{equation*} P\bigl( \{X \in U\} \cap \{Y \in V\} \bigr) = P(\{0\}) = \tfrac{1}{3}, \end{equation*}

whereas

\begin{equation*} P(X \in U) \cdot P(Y \in V) = \tfrac{1}{3} \cdot \tfrac{1}{3} = \tfrac{1}{9} \ne \tfrac{1}{3} . \end{equation*}

Thus the events \(\{X \in U\}\) and \(\{Y \in V\}\) are not independent, and hence \(X\) and \(Y\) are not independent random variables (12.14).

The example works because the additivity of variance is equivalent to the vanishing of \(E(XY) - EX \cdot EY\), and here \(E(XY) = E(X^3) = \frac{1}{3}(-1 + 0 + 1) = 0 = EX \cdot EY\) even though \(Y\) is a (nonconstant) function of \(X\). So independence is genuinely stronger than the conclusion \(\sigma^2(X+Y) = \sigma^2(X) + \sigma^2(Y)\) of 12.22.

Problem (12.12)

Suppose \((\Omega, \mathcal{F}, P)\) and \((\Omega^{\prime}, \mathcal{F}^{\prime}, P^{\prime})\) are probability spaces, \(X\) is a random variable on \(\Omega\), \(Y\) is a random variable on \(\Omega^{\prime}\), and \(\widetilde{X} = \widetilde{Y}\). Prove that \(P_X = P^{\prime}_Y\).

Solution

Here \(P_X\) is the probability measure on \((\mathbb{R}, \mathcal{B})\) defined by \(P_X(B) = P(X \in B)\), and \(P^{\prime}_Y\) is the probability measure on \((\mathbb{R}, \mathcal{B})\) defined by \(P^{\prime}_Y(B) = P^{\prime}(Y \in B)\) (12.27). The hypothesis \(\widetilde{X} = \widetilde{Y}\) says precisely that

\begin{equation*} P_X\bigl( (-\infty, s] \bigr) = P^{\prime}_Y\bigl( (-\infty, s] \bigr) \quad \text{for every } s \in \mathbb{R} . \end{equation*}

We must upgrade this to equality on all Borel sets.

Step 1: an algebra of half-open intervals. Call a subset of \(\mathbb{R}\) a half-open interval if it has one of the forms

\begin{equation*} \emptyset, \quad (a, b], \quad (-\infty, b], \quad (a, \infty), \quad \mathbb{R} \end{equation*}

with \(a, b \in \mathbb{R}\), and let \(\mathcal{A}\) be the collection of all finite unions of half-open intervals.

Two facts about half-open intervals will be used. First, the intersection of two half-open intervals is again a half-open interval: each of the five listed forms can be written as

\begin{equation*} I(\alpha, \beta) = \{ x \in \mathbb{R} : \alpha < x \text{ and } x \le \beta \} \end{equation*}

for suitable \(\alpha \in \mathbb{R} \cup \{-\infty\}\) and \(\beta \in \mathbb{R} \cup \{\infty\}\), and \(I(\alpha_1, \beta_1) \cap I(\alpha_2, \beta_2) = I(\max\{\alpha_1, \alpha_2\}, \min\{\beta_1, \beta_2\})\), which is again one of the five listed forms (it is \(\emptyset\) when \(\max\{\alpha_1, \alpha_2\} \ge \min\{\beta_1, \beta_2\}\)). Second, the complement of a half-open interval is a union of at most two half-open intervals:

\begin{equation*} \mathbb{R} \setminus (a, b] = (-\infty, a] \cup (b, \infty), \quad \mathbb{R} \setminus (-\infty, b] = (b, \infty), \quad \mathbb{R} \setminus (a, \infty) = (-\infty, a], \end{equation*}

together with \(\mathbb{R} \setminus \mathbb{R} = \emptyset\) and \(\mathbb{R} \setminus \emptyset = \mathbb{R}\).

Consequently \(\mathcal{A}\) contains \(\emptyset\), is closed under finite unions by construction, and is closed under complementation: by De Morgan the complement of \(E_1 \cup \cdots \cup E_m\) (each \(E_j\) a half-open interval) is the intersection of the sets \(\mathbb{R} \setminus E_j\), each of which is a union of at most two half-open intervals, and expanding that intersection by distributivity expresses it as a finite union of \(m\)-fold intersections of half-open intervals, hence as a finite union of half-open intervals by the first fact above. Thus \(\mathcal{A}\) is an algebra on \(\mathbb{R}\) in the sense of 5.10.

Step 2: the smallest \(\sigma\)-algebra containing \(\mathcal{A}\) is \(\mathcal{B}\). Every half-open interval is a Borel set, so the smallest \(\sigma\)-algebra containing \(\mathcal{A}\) is contained in \(\mathcal{B}\). Conversely, let \(\mathcal{S}\) denote the smallest \(\sigma\)-algebra containing \(\mathcal{A}\). If \(a < b\) are real numbers then

\begin{equation*} (a, b) = \bigcup_{k=1}^{\infty} \Bigl( a, b - \tfrac{1}{k} \Bigr] \in \mathcal{S}, \end{equation*}

and also \((a, \infty) \in \mathcal{S}\), \((-\infty, b) = \bigcup_{k=1}^{\infty} (-k, b - \frac{1}{k}] \in \mathcal{S}\), and \(\mathbb{R} \in \mathcal{S}\). Thus \(\mathcal{S}\) contains every open interval. Since every open subset of \(\mathbb{R}\) is a countable union of open intervals, \(\mathcal{S}\) contains every open subset of \(\mathbb{R}\), and hence \(\mathcal{S}\) contains the smallest \(\sigma\)-algebra containing the open sets, which is \(\mathcal{B}\). Therefore \(\mathcal{S} = \mathcal{B}\).

Step 3: \(P_X\) and \(P^{\prime}_Y\) agree on \(\mathcal{A}\). Write \(H = \widetilde{X} = \widetilde{Y}\). For \(a < b\) in \(\mathbb{R}\), the sets \((-\infty, a]\) and \((a, b]\) are disjoint with union \((-\infty, b]\), so additivity gives

\begin{equation*} P_X\bigl( (a, b] \bigr) = H(b) - H(a) = P^{\prime}_Y\bigl( (a, b] \bigr) . \end{equation*}

Similarly \(P_X((-\infty, b]) = H(b) = P^{\prime}_Y((-\infty, b])\), and since \(P_X\) and \(P^{\prime}_Y\) are probability measures on \(\mathbb{R}\),

\begin{equation*} P_X\bigl( (a, \infty) \bigr) = 1 - H(a) = P^{\prime}_Y\bigl( (a, \infty) \bigr), \qquad P_X(\mathbb{R}) = 1 = P^{\prime}_Y(\mathbb{R}), \end{equation*}

and \(P_X(\emptyset) = 0 = P^{\prime}_Y(\emptyset)\). So \(P_X\) and \(P^{\prime}_Y\) agree on every half-open interval.

To pass from half-open intervals to all of \(\mathcal{A}\), we show that every element of \(\mathcal{A}\) is a finite union of pairwise disjoint half-open intervals. Suppose \(E = E_1 \cup \cdots \cup E_m\) with each \(E_j\) a half-open interval, written as \(E_j = I(\alpha_j, \beta_j)\) in the notation of Step 1. Let \(c_1 < c_2 < \cdots < c_N\) list the real numbers occurring among \(\alpha_1, \beta_1, \ldots, \alpha_m, \beta_m\) (discarding \(-\infty\) and \(\infty\)). The sets

\begin{equation*} (-\infty, c_1], \quad (c_1, c_2], \quad \ldots, \quad (c_{N-1}, c_N], \quad (c_N, \infty) \end{equation*}

are pairwise disjoint half-open intervals whose union is \(\mathbb{R}\) (if \(N = 0\), take the single set \(\mathbb{R}\)). Each \(\alpha_j\) and each \(\beta_j\) is one of \(c_1, \ldots, c_N, -\infty, \infty\), so each \(E_j = I(\alpha_j, \beta_j)\) is the union of those of the displayed sets that lie between \(\alpha_j\) and \(\beta_j\); hence \(E\) itself is a union of some subcollection of the displayed sets, and that subcollection is pairwise disjoint. Finite additivity of \(P_X\) and of \(P^{\prime}_Y\), together with their agreement on half-open intervals, now gives \(P_X(E) = P^{\prime}_Y(E)\) for every \(E \in \mathcal{A}\).

Step 4: the Monotone Class Theorem. Let

\begin{equation*} \mathcal{M} = \{ B \in \mathcal{B} : P_X(B) = P^{\prime}_Y(B) \} . \end{equation*}

We claim \(\mathcal{M}\) is a monotone class on \(\mathbb{R}\) (5.15). Suppose \(E_1 \subseteq E_2 \subseteq \cdots\) is an increasing sequence in \(\mathcal{M}\). Then by 2.59 applied to each of the two measures,

\begin{equation*} P_X\Bigl( \bigcup_{k=1}^{\infty} E_k \Bigr) = \lim_{k \to \infty} P_X(E_k) = \lim_{k \to \infty} P^{\prime}_Y(E_k) = P^{\prime}_Y\Bigl( \bigcup_{k=1}^{\infty} E_k \Bigr), \end{equation*}

so \(\bigcup_{k=1}^{\infty} E_k \in \mathcal{M}\). Suppose now \(E_1 \supseteq E_2 \supseteq \cdots\) is a decreasing sequence in \(\mathcal{M}\). Because \(P_X(E_1) \le 1 < \infty\) and \(P^{\prime}_Y(E_1) \le 1 < \infty\), the finiteness hypothesis of 2.60 is satisfied, and 2.60 gives

\begin{equation*} P_X\Bigl( \bigcap_{k=1}^{\infty} E_k \Bigr) = \lim_{k \to \infty} P_X(E_k) = \lim_{k \to \infty} P^{\prime}_Y(E_k) = P^{\prime}_Y\Bigl( \bigcap_{k=1}^{\infty} E_k \Bigr), \end{equation*}

so \(\bigcap_{k=1}^{\infty} E_k \in \mathcal{M}\). Thus \(\mathcal{M}\) is a monotone class, and by Step 3 it contains the algebra \(\mathcal{A}\). Hence \(\mathcal{M}\) contains the smallest monotone class containing \(\mathcal{A}\), which by the Monotone Class Theorem (5.17) equals the smallest \(\sigma\)-algebra containing \(\mathcal{A}\), which is \(\mathcal{B}\) by Step 2.

Therefore \(\mathcal{M} = \mathcal{B}\), which says exactly that \(P_X(B) = P^{\prime}_Y(B)\) for every Borel set \(B\). In other words, \(P_X = P^{\prime}_Y\), as desired.

Problem (12.13)

Suppose \(H : \mathbb{R} \to (0, 1)\) is a continuous one-to-one function satisfying conditions (a) through (d) of 12.29. Show that the function \(X : (0, 1) \to \mathbb{R}\) produced in the proof of 12.29 is the inverse function of \(H\).

Solution

Recall that the function produced in the proof of 12.29 is

\begin{equation*} X(\omega) = \sup \{ t \in \mathbb{R} : H(t) < \omega \} \qquad \text{for } \omega \in (0,1) \end{equation*}

(this is 12.30). We must show that \(H\) is a bijection from \(\mathbb{R}\) onto \((0,1)\) and that \(X = H^{-1}\).

Step 1: \(H\) is strictly increasing. Suppose \(s < t\). Condition (a) of 12.29 gives \(H(s) \le H(t)\), and \(H(s) \ne H(t)\) because \(H\) is one-to-one. Hence \(H(s) < H(t)\).

Step 2: \(H\) maps \(\mathbb{R}\) onto \((0,1)\). Let \(u \in (0,1)\). By condition (b), \(\lim_{t \to -\infty} H(t) = 0\), so there exists \(a \in \mathbb{R}\) with \(H(a) < u\). By condition (c), \(\lim_{t \to \infty} H(t) = 1\), so there exists \(b \in \mathbb{R}\) with \(H(b) > u\). Since \(H\) is increasing and \(H(a) < H(b)\), we have \(a < b\). Because \(H\) is continuous on \([a, b]\) and \(H(a) < u < H(b)\), the Intermediate Value Theorem provides \(r \in (a, b)\) with \(H( r) = u\). Since \(H\) takes values in \((0,1)\) by hypothesis, we conclude that \(H\) is a bijection from \(\mathbb{R}\) onto \((0,1)\); in particular the inverse function \(H^{-1} : (0,1) \to \mathbb{R}\) exists.

Step 3: \(X = H^{-1}\). Fix \(\omega \in (0,1)\) and let \(r = H^{-1}(\omega)\), so that \(H( r) = \omega\). We determine the set

\begin{equation*} T_\omega = \{ t \in \mathbb{R} : H(t) < \omega \} \end{equation*}

whose supremum defines \(X(\omega)\).

If \(t < r\), then \(H(t) < H( r) = \omega\) by Step 1, so \(t \in T_\omega\). If \(t \ge r\), then \(H(t) \ge H( r) = \omega\) (using condition (a)), so \(t \notin T_\omega\). Therefore

\begin{equation*} T_\omega = (-\infty, r) . \end{equation*}

This set is nonempty and bounded above by \(r\), and its supremum is \(r\) (every number less than \(r\) belongs to it). Hence

\begin{equation*} X(\omega) = \sup T_\omega = r = H^{-1}(\omega) . \end{equation*}

Since \(\omega \in (0,1)\) was arbitrary, \(X = H^{-1}\), as desired.

As a consistency check, the proof of 12.29 shows that \(\{X \le s\} = \bigl( 0, H(s) \bigr]\) for each \(s \in \mathbb{R}\). With \(X = H^{-1}\) this is immediate: since \(H\) is strictly increasing and onto \((0,1)\), for \(\omega \in (0,1)\) we have \(H^{-1}(\omega) \le s\) if and only if \(\omega = H(H^{-1}(\omega)) \le H(s)\). Because the probability measure on \((0,1)\) in 12.29 is Lebesgue measure, this gives \(\widetilde{X}(s) = \lambda\bigl( (0, H(s)] \bigr) = H(s)\), in agreement with 12.29.

Problem (12.14)

Suppose \((\Omega, \mathcal{F}, P)\) is a probability space and \(X\) is a random variable. Prove that the following are equivalent.

  • \(\widetilde{X}\) is a continuous function on \(\mathbb{R}\).
  • \(\widetilde{X}\) is a uniformly continuous function on \(\mathbb{R}\).
  • \(P(X = t) = 0\) for every \(t \in \mathbb{R}\).
  • \((\widetilde{X} \circ X)^{\sim}(s) = s\) for all \(s \in [0, 1]\).
Solution

Label the four statements (A), (B), (C), (D) in the order listed. Throughout, recall from the first half of the proof of 12.29 that the distribution function \(\widetilde{X}\) is increasing, satisfies \(\lim_{t \to -\infty} \widetilde{X}(t) = 0\) and \(\lim_{t \to \infty} \widetilde{X}(t) = 1\), and is right continuous.

Preliminary computation of the left limits. Fix \(s \in \mathbb{R}\). The sets \(\{X \le s - \frac{1}{k}\}\) for \(k \in \mathbb{Z}^+\) form an increasing sequence of events whose union is \(\{X < s\}\). Hence 2.59 gives

\begin{equation*} \lim_{k \to \infty} \widetilde{X}\Bigl( s - \tfrac{1}{k} \Bigr) = P(X < s) . \end{equation*}

Since \(\widetilde{X}\) is increasing, its left limit at \(s\) exists and equals the limit of any increasing sequence approaching \(s\) from the left; thus

\begin{equation*} \lim_{t \uparrow s} \widetilde{X}(t) = P(X < s) . \end{equation*}

Because \(\{X < s\}\) and \(\{X = s\}\) are disjoint events with union \(\{X \le s\}\), additivity of \(P\) gives

\begin{equation*} \widetilde{X}(s) - \lim_{t \uparrow s} \widetilde{X}(t) = P(X \le s) - P(X < s) = P(X = s) . \end{equation*}

Call this last identity the jump identity; it says that the size of the jump of \(\widetilde{X}\) at \(s\) is exactly \(P(X = s)\).

(A) is equivalent to (C). Since \(\widetilde{X}\) is right continuous at every point, \(\widetilde{X}\) is continuous at \(s\) if and only if \(\lim_{t \uparrow s} \widetilde{X}(t) = \widetilde{X}(s)\), which by the jump identity holds if and only if \(P(X = s) = 0\). Letting \(s\) range over \(\mathbb{R}\) gives the equivalence of (A) and (C).

(B) implies (A). Immediate, since a uniformly continuous function is continuous.

(A) implies (B). Assume \(\widetilde{X}\) is continuous, and let \(\varepsilon > 0\). Since \(\lim_{t \to -\infty} \widetilde{X}(t) = 0\) and \(\lim_{t \to \infty} \widetilde{X}(t) = 1\), there exist \(a < b\) in \(\mathbb{R}\) with

\begin{equation*} \widetilde{X}(a) < \tfrac{\varepsilon}{2} \quad \text{and} \quad \widetilde{X}(b) > 1 - \tfrac{\varepsilon}{2} . \end{equation*}

The function \(\widetilde{X}\) is continuous on the compact interval \([a-1, b+1]\) and hence uniformly continuous there, so there exists \(\delta \in (0, 1)\) such that

\begin{equation*} u, v \in [a-1, b+1], \ |u - v| < \delta \ \implies \ |\widetilde{X}(u) - \widetilde{X}(v)| < \varepsilon . \end{equation*}

Now suppose \(u, v \in \mathbb{R}\) with \(|u - v| < \delta\); say \(u \le v\). There are three cases.

If \(v \le a\), then \(0 \le \widetilde{X}(v) - \widetilde{X}(u) \le \widetilde{X}(a) < \frac{\varepsilon}{2} < \varepsilon\), using that \(\widetilde{X}\) is increasing and nonnegative.

If \(u \ge b\), then \(0 \le \widetilde{X}(v) - \widetilde{X}(u) \le 1 - \widetilde{X}(b) < \frac{\varepsilon}{2} < \varepsilon\), using that \(\widetilde{X} \le 1\) and \(\widetilde{X}\) is increasing.

Otherwise \(v > a\) and \(u < b\). Then \(u > v - \delta > a - \delta > a - 1\) and \(v < u + \delta < b + \delta < b + 1\), so \(u, v \in [a-1, b+1]\) and \(|u - v| < \delta\), giving \(|\widetilde{X}(u) - \widetilde{X}(v)| < \varepsilon\) by the choice of \(\delta\).

In all cases \(|\widetilde{X}(u) - \widetilde{X}(v)| < \varepsilon\). Hence \(\widetilde{X}\) is uniformly continuous on \(\mathbb{R}\).

The random variable \(\widetilde{X} \circ X\). Before proving the remaining implications, note that \(\widetilde{X} : \mathbb{R} \to [0,1]\) is Borel measurable: if \(c \in \mathbb{R}\), then \(\{t \in \mathbb{R} : \widetilde{X}(t) > c\}\) is an interval (if \(\widetilde{X}(t) > c\) and \(t^{\prime} > t\), then \(\widetilde{X}(t^{\prime}) \ge \widetilde{X}(t) > c\) because \(\widetilde{X}\) is increasing), and intervals are Borel sets; thus \(\widetilde{X}\) is Borel measurable by 2.39. Hence \(\widetilde{X} \circ X\) is a random variable on \((\Omega, \mathcal{F})\) by 2.44. Write \(Z = \widetilde{X} \circ X\); its values lie in \([0,1]\), so \(\widetilde{Z}(s) = 0\) for \(s < 0\) and \(\widetilde{Z}(s) = 1\) for \(s \ge 1\). Statement (D) says \(\widetilde{Z}(s) = s\) for all \(s \in [0,1]\).

(A) implies (D). Assume \(\widetilde{X}\) is continuous. Let \(s \in [0,1]\).

Case \(s = 1\): since \(Z \le 1\) everywhere, \(\widetilde{Z}(1) = P(Z \le 1) = 1 = s\).

Case \(s = 0\): since \(Z \ge 0\) everywhere, \(\{Z \le 0\} = \{Z = 0\} = \{X \in N\}\), where

\begin{equation*} N = \{ t \in \mathbb{R} : \widetilde{X}(t) = 0 \} . \end{equation*}

If \(N = \emptyset\) then \(\widetilde{Z}(0) = 0 = s\). Otherwise, \(N\) is nonempty; it is bounded above because \(\lim_{t \to \infty} \widetilde{X}(t) = 1\), and it is closed because \(\widetilde{X}\) is continuous. Let \(u = \sup N\); then \(u \in N\), so \(\widetilde{X}(u) = 0\). Since \(\widetilde{X}\) is increasing, \(N = (-\infty, u]\), and therefore

\begin{equation*} \widetilde{Z}(0) = P(X \le u) = \widetilde{X}(u) = 0 = s . \end{equation*}

Case \(0 < s < 1\): let \(T = \{ t \in \mathbb{R} : \widetilde{X}(t) = s \}\). Because \(\lim_{t \to -\infty} \widetilde{X}(t) = 0 < s\) and \(\lim_{t \to \infty} \widetilde{X}(t) = 1 > s\), there exist \(a < b\) with \(\widetilde{X}(a) < s < \widetilde{X}(b)\); the Intermediate Value Theorem applied to the continuous function \(\widetilde{X}\) on \([a,b]\) shows \(T \ne \emptyset\). Also \(T\) is bounded above by \(b\), and \(T\) is closed by continuity of \(\widetilde{X}\). Let \(b_s = \sup T\); then \(b_s \in T\), so \(\widetilde{X}(b_s) = s\).

We claim \(\{Z \le s\} = \{X \le b_s\}\). Indeed, if \(X(\omega) \le b_s\) then \(Z(\omega) = \widetilde{X}(X(\omega)) \le \widetilde{X}(b_s) = s\) because \(\widetilde{X}\) is increasing. Conversely, if \(X(\omega) > b_s\), put \(t = X(\omega)\); then \(\widetilde{X}(t) \ge \widetilde{X}(b_s) = s\), and \(\widetilde{X}(t) \ne s\) since otherwise \(t \in T\) with \(t > b_s = \sup T\), a contradiction; hence \(Z(\omega) = \widetilde{X}(t) > s\). This proves the claim, and therefore

\begin{equation*} \widetilde{Z}(s) = P(X \le b_s) = \widetilde{X}(b_s) = s . \end{equation*}

Thus (D) holds.

(D) implies (C). We prove the contrapositive. Suppose (C) fails, so there exist \(t_0 \in \mathbb{R}\) and \(c > 0\) with \(P(X = t_0) = c\). Put \(s_0 = \widetilde{X}(t_0)\), so \(s_0 \in [0,1]\), and note that by the jump identity,

\begin{equation*} \lim_{t \uparrow t_0} \widetilde{X}(t) = s_0 - c , \end{equation*}

so in particular \(s_0 - c \ge 0\). If \(t < t_0\) then \(\widetilde{X}(t) \le \lim_{r \uparrow t_0} \widetilde{X}( r) = s_0 - c\), while if \(t \ge t_0\) then \(\widetilde{X}(t) \ge \widetilde{X}(t_0) = s_0\). Hence the range of \(\widetilde{X}\) contains no point of the open interval \((s_0 - c, s_0)\).

Since \(Z = \widetilde{X} \circ X\) takes values in the range of \(\widetilde{X}\), no value of \(Z\) lies in \((s_0 - c, s_0)\). Consequently, for every \(s \in [s_0 - c, s_0)\) we have \(\{Z \le s\} = \{Z \le s_0 - c\}\), and therefore

\begin{equation*} \widetilde{Z}(s) = \widetilde{Z}(s_0 - c) \qquad \text{for all } s \in [s_0 - c, s_0) . \end{equation*}

But \([s_0 - c, s_0) \subseteq [0,1]\), and the two points \(s_0 - c\) and \(s_0 - \frac{c}{2}\) both lie in this interval and are distinct. If (D) held we would get

\begin{equation*} s_0 - \tfrac{c}{2} = \widetilde{Z}\Bigl( s_0 - \tfrac{c}{2} \Bigr) = \widetilde{Z}(s_0 - c) = s_0 - c , \end{equation*}

which is false since \(c > 0\). Hence (D) fails, completing the proof of the contrapositive.

Summary. We have shown (A) is equivalent to (C), (A) is equivalent to (B), (A) implies (D), and (D) implies (C). Since (C) implies (A), all four statements are equivalent.

The content of (D) is the probability integral transform: when \(X\) has no atoms, the random variable \(\widetilde{X} \circ X\) is uniformly distributed on \([0,1]\), which is the converse direction of the construction in the proof of 12.29.

Problem (12.15)

Suppose \(\alpha > 0\) and

\begin{equation*} h(x) = \begin{cases} 0 & \text{if } x < 0, \\ \alpha^2 x e^{-\alpha x} & \text{if } x \geq 0. \end{cases} \end{equation*}

Let \(P = h \, d\lambda\) and let \(X\) be the random variable defined by \(X(x) = x\) for \(x \in \mathbb{R}\).

(a) Verify that \(\int_{-\infty}^{\infty} h \, d\lambda = 1\).

(b) Find a formula for the distribution function \(\widetilde{X}\).

(c) Find a formula (in terms of \(\alpha\)) for \(EX\).

(d) Find a formula (in terms of \(\alpha\)) for \(\sigma(X)\).

Solution

Throughout, \(\lambda\) denotes Lebesgue measure on \(\mathbb{R}\) and \(\mathcal{B}\) denotes the \(\sigma\)-algebra of Borel subsets of \(\mathbb{R}\). Everything follows from one computation, which we isolate.

Lemma. For each integer \(n \geq 0\) and each \(\alpha > 0\),

\begin{equation*} J_n := \int_{[0,\infty)} x^n e^{-\alpha x} \, d\lambda(x) = \frac{n!}{\alpha^{n+1}}. \end{equation*}

Proof of Lemma. The function \(x \mapsto x^n e^{-\alpha x}\) is continuous and nonnegative on \([0, \infty)\), hence Borel measurable. For \(b > 0\) its restriction to \([0,b]\) is continuous, hence Riemann integrable (by 1.11), and by 3.34 its Riemann integral over \([0,b]\) equals its Lebesgue integral over \([0,b]\). Write

\begin{equation*} I_n(b) = \int_0^b x^n e^{-\alpha x} \, dx . \end{equation*}

Then \(I_0(b) = \frac{1 - e^{-\alpha b}}{\alpha}\), and for \(n \geq 1\) integration by parts (with \(u = x^n\), \(dv = e^{-\alpha x} dx\)) gives

\begin{equation*} I_n(b) = -\frac{b^n e^{-\alpha b}}{\alpha} + \frac{n}{\alpha} I_{n-1}(b). \end{equation*}

Now for \(k \in \mathbb{Z}^+\) define \(f_k \colon \mathbb{R} \to [0,\infty)\) by \(f_k(x) = x^n e^{-\alpha x} \chi_{[0,k]}(x)\), and define \(f(x) = x^n e^{-\alpha x} \chi_{[0,\infty)}(x)\). Then \(0 \leq f_1 \leq f_2 \leq \cdots\), and \(f_k(x) \to f(x)\) for every \(x \in \mathbb{R}\) (given \(x\), we have \(f_k(x) = f(x)\) as soon as \(k \geq x\)). Since \(\int f_k \, d\lambda = I_n(k)\), the Monotone Convergence Theorem (3.11) gives

\begin{equation*} J_n = \int f \, d\lambda = \lim_{k \to \infty} I_n(k). \end{equation*}

Since \(b^n e^{-\alpha b} \to 0\) as \(b \to \infty\), and since \(b \mapsto I_n(b)\) is increasing (its integrand is nonnegative) so that \(\lim_{b \to \infty} I_n(b) = \lim_{k \to \infty} I_n(k) = J_n\), letting \(b \to \infty\) in the two displayed recursions yields \(J_0 = \frac{1}{\alpha}\) and \(J_n = \frac{n}{\alpha} J_{n-1}\) for \(n \geq 1\). Induction now gives \(J_n = n!/\alpha^{n+1}\), completing the proof of the Lemma.

Part (a). The function \(h\) is Borel measurable (it is continuous on \((-\infty,0)\) and on \([0,\infty)\), and each of these sets is Borel) and \(h \geq 0\). Since \(h\) vanishes on \((-\infty, 0)\), the Lemma with \(n = 1\) gives

\begin{equation*} \int_{-\infty}^{\infty} h \, d\lambda = \alpha^2 \int_{[0,\infty)} x e^{-\alpha x} \, d\lambda(x) = \alpha^2 J_1 = \alpha^2 \cdot \frac{1!}{\alpha^2} = 1 . \end{equation*}

In particular \(h \in L^1(\mathbb{R})\), so the hypotheses of 12.33 are satisfied by \(h\), by the probability measure \(P\) on \((\mathbb{R}, \mathcal{B})\) defined by \(P(B) = \int_B h \, d\lambda\), and by the random variable \(X(x) = x\).

Part (b). By 12.33, \(h\) is the density function of \(X\); thus

\begin{equation*} \widetilde{X}(s) = P(X \leq s) = \int_{-\infty}^{s} h \, d\lambda \end{equation*}

for all \(s \in \mathbb{R}\). If \(s < 0\) then \(h = 0\) on \((-\infty, s]\), so \(\widetilde{X}(s) = 0\). If \(s \geq 0\), then on \([0,s]\) the integrand is continuous, so (by 3.34 again) the Lebesgue integral equals the Riemann integral, and since

\begin{equation*} \frac{d}{dx}\Bigl[ -(1 + \alpha x) e^{-\alpha x} \Bigr] = -\alpha e^{-\alpha x} + \alpha(1 + \alpha x) e^{-\alpha x} = \alpha^2 x e^{-\alpha x}, \end{equation*}

the Fundamental Theorem of Calculus gives

\begin{equation*} \int_{-\infty}^{s} h \, d\lambda = \int_0^s \alpha^2 x e^{-\alpha x} \, dx = \Bigl[ -(1 + \alpha x) e^{-\alpha x} \Bigr]_0^s = 1 - (1 + \alpha s) e^{-\alpha s}. \end{equation*}

Hence

\begin{equation*} \widetilde{X}(s) = \begin{cases} 0 & \text{if } s < 0, \\ 1 - (1 + \alpha s) e^{-\alpha s} & \text{if } s \geq 0. \end{cases} \end{equation*}

(As a check, \(\widetilde{X}\) is continuous, increasing, \(\widetilde{X}(s) \to 0\) as \(s \to -\infty\), and \(\widetilde{X}(s) \to 1\) as \(s \to \infty\), which is consistent with (a).)

Part (c). To apply the mean formula in 12.33 we must first verify that \(X \in L^1(P)\), and for that we need the change-of-measure identity

\begin{equation*} \int_{\mathbb{R}} f \, dP = \int_{\mathbb{R}} f h \, d\lambda \end{equation*}

for every Borel measurable \(f \colon \mathbb{R} \to [0,\infty]\). (The book records this identity for \(f \in L^1(P)\) in the proof of 12.33, citing Exercise 5 in Section 9A; the nonnegative case, which is what we need, follows from the definition of \(P\) as follows.) If \(f = \chi_B\) with \(B \in \mathcal{B}\), the identity is exactly the definition \(P(B) = \int_B h \, d\lambda\). If \(f = c \chi_B\) with \(c \in [0,\infty)\) and \(B \in \mathcal{B}\), then since \(\int_B h \, d\lambda\) is finite (it is at most \(\int_{\mathbb{R}} h \, d\lambda = 1\), by 3.8) we may use 3.20 twice to get \(\int c \chi_B \, dP = c P(B) = c \int_B h \, d\lambda = \int c \chi_B h \, d\lambda\). Since every nonnegative simple Borel function is a finite sum of such \(c \chi_B\), additivity of the integral of nonnegative functions (3.16) shows that the identity holds for every nonnegative simple Borel function. For general Borel \(f \geq 0\), take by 2.89 simple Borel functions \(0 \leq f_1 \leq f_2 \leq \cdots\) with \(f_k \to f\) pointwise; then \(0 \leq f_1 h \leq f_2 h \leq \cdots\) and \(f_k h \to f h\) pointwise, so two applications of the Monotone Convergence Theorem (3.11), one for \(P\) and one for \(\lambda\), give

\begin{equation*} \int_{\mathbb{R}} f \, dP = \lim_{k \to \infty} \int_{\mathbb{R}} f_k \, dP = \lim_{k \to \infty} \int_{\mathbb{R}} f_k h \, d\lambda = \int_{\mathbb{R}} f h \, d\lambda . \end{equation*}

Applying this with \(f = |X|\), the Lemma with \(n = 2\) gives

\begin{equation*} \int_{\mathbb{R}} |X| \, dP = \int_{-\infty}^{\infty} |x| h(x) \, d\lambda(x) = \alpha^2 J_2 = \alpha^2 \cdot \frac{2!}{\alpha^3} = \frac{2}{\alpha} < \infty . \end{equation*}

Therefore, by the formula for the mean in 12.33,

\begin{equation*} EX = \int_{-\infty}^{\infty} x h(x) \, d\lambda(x) = \alpha^2 J_2 = \frac{2}{\alpha}. \end{equation*}

Part (d). Applying the change-of-measure identity of part (c) with \(f = |X|^2\), and then the Lemma with \(n = 3\), we get

\begin{equation*} \int_{\mathbb{R}} |X|^2 \, dP = \int_{-\infty}^{\infty} x^2 h(x) \, d\lambda(x) = \alpha^2 J_3 = \alpha^2 \cdot \frac{3!}{\alpha^4} = \frac{6}{\alpha^2} < \infty, \end{equation*}

so \(X \in L^2(P)\). Hence the variance formula in 12.33 applies:

\begin{equation*} \sigma^2(X) = \int_{-\infty}^{\infty} x^2 h(x) \, d\lambda(x) - \left( \int_{-\infty}^{\infty} x h(x) \, d\lambda(x) \right)^{2} = \frac{6}{\alpha^2} - \frac{4}{\alpha^2} = \frac{2}{\alpha^2}. \end{equation*}

Taking square roots,

\begin{equation*} \sigma(X) = \frac{\sqrt{2}}{\alpha}. \end{equation*}

(Sanity check, not used above: \(h\) is the density of a sum of two independent exponential random variables each with parameter \(\alpha\), and the second bullet point of 12.34 gives mean \(1/\alpha\) and variance \(1/\alpha^2\) for each summand; adding means and, by 12.22, adding variances predicts \(2/\alpha\) and \(2/\alpha^2\), in agreement with the answers found above.)

Problem (12.16)

Suppose \(\mathcal{B}\) is the \(\sigma\)-algebra of Borel subsets of \([0,1)\) and \(P\) is Lebesgue measure on \(\bigl( [0,1), \mathcal{B} \bigr)\). Let \(\{e_k\}_{k \in \mathbb{Z}^+}\) be the family of functions defined by the fourth bullet point of Example 8.51 (notice that \(k = 0\) is excluded). Show that the family \(\{e_k\}_{k \in \mathbb{Z}^+}\) is an i.i.d.

Solution

Recall the fourth bullet point of 8.51: for a nonnegative integer \(k\), the Rademacher function \(e_k \colon [0,1) \to \mathbb{R}\) is defined by

\begin{equation*} e_k(x) = \begin{cases} 1 & \text{if } x \in \bigl[ \tfrac{n-1}{2^k}, \tfrac{n}{2^k} \bigr) \text{ for some odd integer } n, \\ -1 & \text{if } x \in \bigl[ \tfrac{n-1}{2^k}, \tfrac{n}{2^k} \bigr) \text{ for some even integer } n. \end{cases} \end{equation*}

For \(k \in \mathbb{Z}^+\) and \(n \in \{1, \dots, 2^k\}\) put

\begin{equation*} I_{k,n} = \Bigl[ \tfrac{n-1}{2^k}, \tfrac{n}{2^k} \Bigr). \end{equation*}

These \(2^k\) intervals are pairwise disjoint, their union is \([0,1)\), and each has \(P(I_{k,n}) = 2^{-k}\); the intervals \([\frac{n-1}{2^k},\frac{n}{2^k})\) for other integers \(n\) are disjoint from \([0,1)\). Thus \(e_k\) is a well-defined function on \([0,1)\) taking only the values \(1\) and \(-1\), and it is a finite linear combination of characteristic functions of intervals, hence Borel measurable. So each \(e_k\) is a random variable on \(([0,1), \mathcal{B})\). Note also that \(P([0,1)) = 1\), so \(P\) is indeed a probability measure and \(([0,1), \mathcal{B}, P)\) is a probability space.

Step 1: each \(e_k\) takes the values \(\pm 1\) with probability \(\tfrac12\).

Fix \(k \in \mathbb{Z}^+\). Exactly \(2^{k-1}\) of the integers \(n \in \{1,\dots,2^k\}\) are odd and exactly \(2^{k-1}\) are even (here \(k \geq 1\) is used). Since \(\{e_k = 1\}\) is the union of the \(I_{k,n}\) with \(n\) odd and \(\{e_k = -1\}\) is the union of those with \(n\) even, countable additivity gives

\begin{equation*} P(e_k = 1) = 2^{k-1} \cdot 2^{-k} = \tfrac12, \qquad P(e_k = -1) = 2^{k-1} \cdot 2^{-k} = \tfrac12 . \end{equation*}

Step 2: the family is identically distributed.

By Step 1, for every \(k \in \mathbb{Z}^+\) the distribution function of \(e_k\) is

\begin{equation*} \widetilde{e_k}(s) = P(e_k \leq s) = \begin{cases} 0 & \text{if } s < -1, \\ \tfrac12 & \text{if } -1 \leq s < 1, \\ 1 & \text{if } s \geq 1. \end{cases} \end{equation*}

This formula does not involve \(k\), so \(P(e_j \leq s) = P(e_k \leq s)\) for all \(j,k \in \mathbb{Z}^+\) and all \(s \in \mathbb{R}\). Hence \(\{e_k\}_{k \in \mathbb{Z}^+}\) is identically distributed (12.35). (This is exactly where \(k = 0\) must be excluded: the interval \([0,1)\) is \(I_{0,1}\) with \(n = 1\) odd, so \(e_0\) is identically \(1\) and \(\widetilde{e_0}\) differs from the function displayed above.)

Step 3: \(e_k\) reads off the \(k\)-th binary digit.

Rather than invoke the theory of binary expansions, we define the binary digits directly and prove what we need. For \(x \in [0,1)\) and \(j \in \mathbb{Z}^+\) set

\begin{equation*} d_j(x) = \lfloor 2^j x \rfloor - 2 \lfloor 2^{j-1} x \rfloor . \end{equation*}

First, \(d_j(x) \in \{0,1\}\): writing \(a = 2^{j-1}x\), we have \(\lfloor 2a \rfloor - 2\lfloor a \rfloor = \lfloor 2(a - \lfloor a \rfloor) \rfloor\), and \(a - \lfloor a \rfloor \in [0,1)\), so this lies in \(\{0,1\}\). Second, since \(\lfloor x \rfloor = 0\) for \(x \in [0,1)\), the sum \(\sum_{j=1}^k d_j(x) 2^{k-j}\) telescopes:

\begin{equation*} \sum_{j=1}^{k} d_j(x) 2^{k-j} = \sum_{j=1}^{k} \Bigl( \lfloor 2^j x \rfloor 2^{k-j} - \lfloor 2^{j-1} x \rfloor 2^{k-j+1} \Bigr) = \lfloor 2^k x \rfloor - 2^k \lfloor x \rfloor = \lfloor 2^k x \rfloor . \end{equation*}

Thus

\begin{equation*} \lfloor 2^k x \rfloor = \sum_{j=1}^{k} d_j(x) 2^{k-j}. \end{equation*}

On the other hand, \(x \in I_{k,n}\) if and only if \(n - 1 = \lfloor 2^k x \rfloor\). So the unique \(n \in \{1,\dots,2^k\}\) with \(x \in I_{k,n}\) is

\begin{equation*} n = 1 + \sum_{j=1}^{k} d_j(x) 2^{k-j}. \end{equation*}

Every term of this sum with \(j < k\) is even, so \(n \equiv 1 + d_k(x) \pmod 2\). Therefore \(n\) is odd if and only if \(d_k(x) = 0\), which gives

\begin{equation*} e_k(x) = 1 \iff d_k(x) = 0, \qquad e_k(x) = -1 \iff d_k(x) = 1 . \end{equation*}

The displayed formula for \(n\) also shows the following: for each \(N \in \mathbb{Z}^+\) the map \(x \mapsto (d_1(x), \dots, d_N(x))\) is constant on each \(I_{N,n}\), and since \((a_1,\dots,a_N) \mapsto 1 + \sum_{j=1}^N a_j 2^{N-j}\) is a bijection from \(\{0,1\}^N\) onto \(\{1,\dots,2^N\}\), each element of \(\{0,1\}^N\) is attained on exactly one of the \(2^N\) intervals \(I_{N,1}, \dots, I_{N,2^N}\).

Step 4: the probability of prescribing finitely many values.

Let \(m \in \mathbb{Z}^+\), let \(k_1 < k_2 < \cdots < k_m\) be elements of \(\mathbb{Z}^+\), and let \(\varepsilon_1, \dots, \varepsilon_m \in \{-1,1\}\). Put \(\delta_i = \frac{1 - \varepsilon_i}{2} \in \{0,1\}\) and set \(N = k_m\). By Step 3,

\begin{equation*} E := \bigcap_{i=1}^{m} \{ e_{k_i} = \varepsilon_i \} = \{ x \in [0,1) : d_{k_i}(x) = \delta_i \text{ for } i = 1, \dots, m \}. \end{equation*}

By the last paragraph of Step 3, \(E\) is the union of exactly those intervals \(I_{N,n}\) whose associated digit string \((a_1,\dots,a_N) \in \{0,1\}^N\) satisfies \(a_{k_i} = \delta_i\) for \(i = 1,\dots,m\). Since \(k_1, \dots, k_m\) are \(m\) distinct elements of \(\{1,\dots,N\}\), exactly \(2^{N-m}\) of the \(2^N\) digit strings satisfy those constraints. As the intervals \(I_{N,n}\) are disjoint and each has measure \(2^{-N}\),

\begin{equation*} P(E) = 2^{N-m} \cdot 2^{-N} = 2^{-m} = \prod_{i=1}^{m} P(e_{k_i} = \varepsilon_i), \end{equation*}

the last equality by Step 1.

Step 5: the family is independent.

Let \(\{U_k\}_{k \in \mathbb{Z}^+}\) be any family of Borel subsets of \(\mathbb{R}\); we must show that the family of events \(\{ e_k \in U_k \}_{k \in \mathbb{Z}^+}\) is independent in the sense of 12.7. So let \(k_1, \dots, k_m\) be distinct elements of \(\mathbb{Z}^+\); relabeling if necessary (both sides of the identity to be proved are unchanged by permuting the indices) we may assume \(k_1 < \cdots < k_m\). Write \(V_i = U_{k_i} \cap \{-1, 1\}\). Because \(e_{k_i}\) takes only the values \(\pm 1\),

\begin{equation*} \{ e_{k_i} \in U_{k_i} \} = \bigcup_{\varepsilon \in V_i} \{ e_{k_i} = \varepsilon \}, \qquad \text{so} \qquad P(e_{k_i} \in U_{k_i}) = \tfrac{1}{2} \, |V_i| \end{equation*}

by Step 1 (here \(|V_i| \in \{0,1,2\}\) denotes the number of elements of \(V_i\), and the union is disjoint). Distributing the intersection over these unions gives

\begin{equation*} \bigcap_{i=1}^{m} \{ e_{k_i} \in U_{k_i} \} = \bigcup_{(\varepsilon_1,\dots,\varepsilon_m) \in V_1 \times \cdots \times V_m} \; \bigcap_{i=1}^{m} \{ e_{k_i} = \varepsilon_i \}, \end{equation*}

and the sets in this union are pairwise disjoint, because distinct tuples \((\varepsilon_1,\dots,\varepsilon_m)\) prescribe different values for some \(e_{k_i}\). There are \(|V_1| \cdots |V_m|\) such tuples, and by Step 4 each corresponding intersection has probability \(2^{-m}\). Hence

\begin{equation*} P\Bigl( \bigcap_{i=1}^{m} \{ e_{k_i} \in U_{k_i} \} \Bigr) = |V_1| \cdots |V_m| \cdot 2^{-m} = \prod_{i=1}^{m} \frac{|V_i|}{2} = \prod_{i=1}^{m} P(e_{k_i} \in U_{k_i}). \end{equation*}

(If some \(V_i\) is empty, both sides are \(0\), consistent with the formula.) Thus \(\{ e_k \in U_k \}_{k \in \mathbb{Z}^+}\) is an independent family of events for every family \(\{U_k\}_{k \in \mathbb{Z}^+}\) of Borel subsets of \(\mathbb{R}\), which by 12.14 says precisely that \(\{e_k\}_{k \in \mathbb{Z}^+}\) is an independent family of random variables.

Combining Step 2 and Step 5, the family \(\{e_k\}_{k \in \mathbb{Z}^+}\) is independent and identically distributed, i.e., it is an i.i.d. family (12.35).

Problem (12.17)

Suppose \(\mathcal{B}\) is the \(\sigma\)-algebra of Borel subsets of \((-\pi, \pi]\) and \(P\) is Lebesgue measure on \(\bigl( (-\pi,\pi], \mathcal{B} \bigr)\) divided by \(2\pi\). Let \(\{e_k\}_{k \in \mathbb{Z} \setminus \{0\}}\) be the family of trigonometric functions defined by the third bullet point of Example 8.51 (notice that \(k = 0\) is excluded).

(a) Show that \(\{e_k\}_{k \in \mathbb{Z} \setminus \{0\}}\) is not an independent family of random variables.

(b) Show that \(\{e_k\}_{k \in \mathbb{Z} \setminus \{0\}}\) is an identically distributed family.

Solution

Recall the third bullet point of 8.51: for \(k \in \mathbb{Z}\) the function \(e_k \colon (-\pi,\pi] \to \mathbb{R}\) is defined by

\begin{equation*} e_k(t) = \begin{cases} \frac{1}{\sqrt{\pi}} \sin(kt) & \text{if } k > 0, \\ \frac{1}{\sqrt{2\pi}} & \text{if } k = 0, \\ \frac{1}{\sqrt{\pi}} \cos(kt) & \text{if } k < 0. \end{cases} \end{equation*}

Each \(e_k\) is continuous, hence Borel measurable, so each \(e_k\) is a random variable on \(((-\pi,\pi], \mathcal{B})\). Throughout, \(\lambda\) denotes Lebesgue measure on \(\mathbb{R}\), so that \(P(A) = \frac{\lambda(A)}{2\pi}\) for \(A \in \mathcal{B}\); note \(P\bigl( (-\pi,\pi] \bigr) = 1\), so \(P\) is indeed a probability measure.

Part (a). It suffices to produce two indices in \(\mathbb{Z} \setminus \{0\}\) and two Borel subsets of \(\mathbb{R}\) for which the product formula of 12.7 fails; we use \(k = 1\) and \(k = -1\), so that

\begin{equation*} e_1(t) = \frac{\sin t}{\sqrt{\pi}}, \qquad e_{-1}(t) = \frac{\cos(-t)}{\sqrt{\pi}} = \frac{\cos t}{\sqrt{\pi}} . \end{equation*}

The point is the Pythagorean identity: \(e_1(t)^2 + e_{-1}(t)^2 = \frac{1}{\pi}\) for every \(t\), so these two random variables are very far from independent. Let

\begin{equation*} U = \Bigl\{ x \in \mathbb{R} : x^2 > \tfrac{1}{2\pi} \Bigr\}, \end{equation*}

a Borel (indeed open) subset of \(\mathbb{R}\), and set \(U_1 = U_{-1} = U\) and \(U_k = \mathbb{R}\) for all other \(k \in \mathbb{Z} \setminus \{0\}\).

First, \(\{ e_1 \in U \} \cap \{ e_{-1} \in U \} = \varnothing\): if \(t\) belonged to both sets we would have

\begin{equation*} \frac{1}{\pi} = e_1(t)^2 + e_{-1}(t)^2 > \frac{1}{2\pi} + \frac{1}{2\pi} = \frac{1}{\pi}, \end{equation*}

a contradiction. Hence

\begin{equation*} P\bigl( \{ e_1 \in U \} \cap \{ e_{-1} \in U \} \bigr) = 0 . \end{equation*}

Next we compute the two individual probabilities. Since \(e_1(t)^2 > \frac{1}{2\pi}\) is equivalent to \(|\sin t| > \frac{1}{\sqrt{2}}\), we have

\begin{equation*} \{ e_1 \in U \} = \Bigl( -\tfrac{3\pi}{4}, -\tfrac{\pi}{4} \Bigr) \cup \Bigl( \tfrac{\pi}{4}, \tfrac{3\pi}{4} \Bigr), \end{equation*}

which has Lebesgue measure \(\frac{\pi}{2} + \frac{\pi}{2} = \pi\); thus \(P(e_1 \in U) = \frac{\pi}{2\pi} = \frac12\). Similarly \(e_{-1}(t)^2 > \frac{1}{2\pi}\) is equivalent to \(|\cos t| > \frac{1}{\sqrt{2}}\), so

\begin{equation*} \{ e_{-1} \in U \} = \Bigl( -\pi, -\tfrac{3\pi}{4} \Bigr) \cup \Bigl( -\tfrac{\pi}{4}, \tfrac{\pi}{4} \Bigr) \cup \Bigl( \tfrac{3\pi}{4}, \pi \Bigr], \end{equation*}

which has Lebesgue measure \(\frac{\pi}{4} + \frac{\pi}{2} + \frac{\pi}{4} = \pi\); thus \(P(e_{-1} \in U) = \frac12\). Therefore

\begin{equation*} P\bigl( \{ e_1 \in U_1 \} \cap \{ e_{-1} \in U_{-1} \} \bigr) = 0 \neq \tfrac14 = P(e_1 \in U_1) \cdot P(e_{-1} \in U_{-1}). \end{equation*}

So the family of events \(\{ e_k \in U_k \}_{k \in \mathbb{Z} \setminus \{0\}}\) is not independent (12.7), and hence, by 12.14, the family \(\{e_k\}_{k \in \mathbb{Z} \setminus \{0\}}\) is not an independent family of random variables. (In fact this shows more: \(e_1\) and \(e_{-1}\) are not even pairwise independent.)

Part (b). We must show that \(P(e_j \leq s) = P(e_k \leq s)\) for all \(j, k \in \mathbb{Z} \setminus \{0\}\) and all \(s \in \mathbb{R}\). Call a set \(S \subseteq \mathbb{R}\) periodic if \(S + 2\pi = S\). We first record two elementary lemmas.

Lemma A. If \(S \subseteq \mathbb{R}\) is a periodic Borel set and \(a \in \mathbb{R}\), then \(\lambda\bigl( S \cap (a, a + 2\pi] \bigr) = \lambda\bigl( S \cap (-\pi, \pi] \bigr)\).

Proof of Lemma A. Translating \(a\) by an integer multiple of \(2\pi\) changes neither side: if \(a^{\prime} = a + 2\pi m\) with \(m \in \mathbb{Z}\), then \(S \cap (a^{\prime}, a^{\prime}+2\pi] = \bigl( S \cap (a,a+2\pi] \bigr) + 2\pi m\) because \(S\) is periodic, and Lebesgue measure is translation invariant (2.7). So we may assume \(-\pi < a \leq \pi\). Then

\begin{equation*} S \cap (a, a+2\pi] = \bigl( S \cap (a, \pi] \bigr) \cup \bigl( S \cap (\pi, a + 2\pi] \bigr), \end{equation*}

a disjoint union. Since \((\pi, a+2\pi] = (-\pi, a] + 2\pi\) and \(S\) is periodic, \(S \cap (\pi, a+2\pi] = \bigl( S \cap (-\pi, a] \bigr) + 2\pi\), so these two sets have the same measure by 2.7. Hence

\begin{equation*} \lambda\bigl( S \cap (a, a+2\pi] \bigr) = \lambda\bigl( S \cap (a,\pi] \bigr) + \lambda\bigl( S \cap (-\pi, a] \bigr) = \lambda\bigl( S \cap (-\pi,\pi] \bigr), \end{equation*}

as claimed.

Lemma B. If \(S \subseteq \mathbb{R}\) is a periodic Borel set and \(m \in \mathbb{Z}^+\), then

\begin{equation*} \lambda\bigl( \{ t \in (-\pi, \pi] : mt \in S \} \bigr) = \lambda\bigl( S \cap (-\pi, \pi] \bigr). \end{equation*}

Proof of Lemma B. For \(t \in (-\pi,\pi]\) we have \(mt \in (-m\pi, m\pi]\), so

\begin{equation*} \{ t \in (-\pi,\pi] : mt \in S \} = \tfrac{1}{m} \bigl( S \cap (-m\pi, m\pi] \bigr). \end{equation*}

By 5.41 (with \(n = 1\) and dilation factor \(\frac1m > 0\)) this set is Borel and its measure equals \(\frac{1}{m} \lambda\bigl( S \cap (-m\pi, m\pi] \bigr)\). Now \((-m\pi, m\pi]\) is the disjoint union of the \(m\) intervals \(\bigl( -m\pi + 2\pi j, \, -m\pi + 2\pi(j+1) \bigr]\) for \(j = 0, 1, \dots, m-1\), each of length \(2\pi\), so Lemma A gives

\begin{equation*} \lambda\bigl( S \cap (-m\pi, m\pi] \bigr) = \sum_{j=0}^{m-1} \lambda\bigl( S \cap (-m\pi + 2\pi j, -m\pi + 2\pi (j+1)] \bigr) = m \, \lambda\bigl( S \cap (-\pi,\pi] \bigr). \end{equation*}

Dividing by \(m\) proves Lemma B.

Now fix \(s \in \mathbb{R}\) and set

\begin{equation*} S = \{ u \in \mathbb{R} : \sin u \leq s \sqrt{\pi} \}, \qquad T = \{ u \in \mathbb{R} : \cos u \leq s \sqrt{\pi} \}. \end{equation*}

Both are Borel (preimages of the closed set \((-\infty, s\sqrt{\pi}\,]\) under continuous functions) and both are periodic, since \(\sin\) and \(\cos\) have period \(2\pi\).

Suppose first that \(k > 0\). Then \(e_k(t) \leq s\) if and only if \(\sin(kt) \leq s\sqrt{\pi}\), i.e., if and only if \(kt \in S\). Hence by Lemma B (with \(m = k\)),

\begin{equation*} P(e_k \leq s) = \frac{1}{2\pi} \lambda\bigl( \{ t \in (-\pi,\pi] : kt \in S \} \bigr) = \frac{1}{2\pi} \lambda\bigl( S \cap (-\pi,\pi] \bigr). \end{equation*}

The right side does not depend on \(k\).

Now suppose \(k < 0\), and put \(m = -k \in \mathbb{Z}^+\). Since cosine is an even function, \(e_k(t) = \frac{1}{\sqrt{\pi}} \cos(kt) = \frac{1}{\sqrt{\pi}} \cos(mt)\). Thus \(e_k(t) \leq s\) if and only if \(mt \in T\), so Lemma B gives

\begin{equation*} P(e_k \leq s) = \frac{1}{2\pi} \lambda\bigl( T \cap (-\pi,\pi] \bigr). \end{equation*}

It remains to check that \(\lambda\bigl( T \cap (-\pi,\pi] \bigr) = \lambda\bigl( S \cap (-\pi,\pi] \bigr)\). Because \(\cos u = \sin\bigl( u + \frac{\pi}{2} \bigr)\), we have \(u \in T\) if and only if \(u + \frac{\pi}{2} \in S\); that is, \(T = S - \frac{\pi}{2}\). Consequently

\begin{equation*} T \cap (-\pi, \pi] = \Bigl( S \cap \Bigl( -\tfrac{\pi}{2}, \tfrac{3\pi}{2} \Bigr] \Bigr) - \tfrac{\pi}{2}, \end{equation*}

so by translation invariance (2.7) and then Lemma A applied to the length-\(2\pi\) interval \(\bigl( -\frac{\pi}{2}, \frac{3\pi}{2} \bigr]\),

\begin{equation*} \lambda\bigl( T \cap (-\pi,\pi] \bigr) = \lambda\Bigl( S \cap \Bigl( -\tfrac{\pi}{2}, \tfrac{3\pi}{2} \Bigr] \Bigr) = \lambda\bigl( S \cap (-\pi,\pi] \bigr). \end{equation*}

Therefore, for every \(k \in \mathbb{Z} \setminus \{0\}\) and every \(s \in \mathbb{R}\),

\begin{equation*} \widetilde{e_k}(s) = P(e_k \leq s) = \frac{1}{2\pi} \lambda\bigl( \{ t \in (-\pi,\pi] : \sin t \leq s \sqrt{\pi} \} \bigr), \end{equation*}

a quantity independent of \(k\). By 12.35, the family \(\{e_k\}_{k \in \mathbb{Z}\setminus\{0\}}\) is identically distributed.

(Note that \(k = 0\) must be excluded: \(e_0\) is the constant \(\frac{1}{\sqrt{2\pi}}\), whose distribution function jumps from \(0\) to \(1\) at \(\frac{1}{\sqrt{2\pi}}\), whereas the common distribution function above is continuous. Note also that parts (a) and (b) together give an example showing that an identically distributed family need not be independent.)